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Identification, Estimation and Inference Based on Structural Error Projection

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titlepage\begin{center} Identification, Estimation and Inference Based on \\Structural Error Projection\footnote{An earlier version of this paper titled “Semiparametric Instrumental Variables Method" has been discussed and presented at several places. We would like to thank them for their constructive comments and suggestions, particularly to Otavio Bartalotti, Debopam Bhattacharya, Giuseppe Cavalieve, Xiaohong Chen, Benjamin Deaner, Firmin Doko Tchatoka, James Duffy, Yanqin Fan, David Frazier, Silvia Gonçalves, Yongmiao Hong, Arturas Juodis, Toru Kitagawa, Frank Kleibergen, Tatiana Komarova, Dennis Kristensen, Sophocles Mavroedis, Akanksha Negi, Whitney Newey, Didier Nibbering, Taisuke Otsu, Hashem Pesaran, Peter Phillips, Taiga Saito, Richard Smith, Liangjun Su, Yike Wang, Martin Weidner, Ruofan Xu, Jun Yu, Lina Zhang, and Xueyan Zhao. Dong acknowledges from the National Natural Science Foundation of China under Grant Numbers: 72473156 and 72073143. Gao, Linton and Peng would like to thank the Australian Research Council Discovery Grants Program for its financial support under Grant Number: DP250100063. Thanks also go to Ruofan Xu for her assistance in computing. An online supplementary document and Matlab code are both available at \url{https://github.com/pengbin430/SIV_2026/tree/main.}} {\sc Chaohua Dong$^{\ast}$, Jiti Gao$^{\dag}$, Oliver Linton$^{\star}$ and Bin Peng$^{\dag}$} { $^{\ast}$Zhongnan University of Economics and Law, $^{\dag}$Monash University and $^{\star}$University of Cambridge} \end{center} \begin{abstract} This paper proposes to project and expand the conditional mean function of the structural error given the regressors in an endogenous regression under consideration. As the projection process is semiparametric, we define this procedure as a semiparametric projection (SP) method to address endogeneity in regression models by internally constructed instrumental variables. The SP method is applicable to many classes of regression models associated with endogeneity, such as linear, nonlinear, and non-- and semi--parametric models, and provides a simple and computationally tractable alternative to conventional instrumental variable approaches available from the existing literature. This paper establishes identification conditions and derives the asymptotic properties of the resulting estimators. It then proposes a simple LASSO selection method to examine the finite--sample performance of both the proposed method and the established theory by simulated and real data examples. \end{abstract} Some key words: Causal effect; Control Function, Endogenous regressors, Instrumental Variable; Projection; Structural Equation JEL subject classification: C14, C26, G51

Introduction

Identifying structural parameters when regressors are correlated with unobserved disturbances is a central problem in econometrics. Instrumental variable (IV) methods provide a widely used solution, but they rely on the availability of valid instruments that are correlated with the endogenous regressors but uncorrelated with the structural errors. In many empirical settings such instruments are difficult to obtain or justify. This paper proposes a semiparametric instrumental variable approach that constructs instruments directly from the observed regressors. The key observation is that endogeneity arises from the component of the structural error that is correlated with the regressors. By projecting the structural equation onto the sigma--field generated by the regressors, we decompose the structural error into a conditional expectation component and an orthogonal residual. This decomposition yields an exogenous regression representation that allows us to construct valid instruments as functions of the observed covariates.

Suppose that we have a linear structural model

equation[equation omitted — 162 chars of source]

where $\beta _{0}$ are the parameters of interest and $\mathbf{x}$ is the observed $d$-dimensional vector of variables related to the outcome $y$. The main issue we face regarding identification and estimation is that ${\mathbb{ E}}[\mathbf{x}\varepsilon ]\neq 0.$ We write $\varepsilon ={\mathbb{E}} [\varepsilon |\mathbf{x}]+e,$ where $m(\mathbf{x)}={\mathbb{E}}[\varepsilon | \mathbf{x}]$ is of unknown functional form and ${\mathbb{E}}[e|\mathbf{x}]=0$ . We suppose that $m\in S,$ a subspace of the Hilbert space of square integrable functions of $\mathbf{x}$. Letting $\mathbf{P}_{\mathcal{S}}$ denote the projection operator on to the subspace $\mathcal{S}$, we define instruments $\mathbf{z}=\mathbf{x}- \mathbf{P}_{\mathcal{S}} \, \mathbf{x}$, and then apply the IV method to estimate the structural parameters. Because $ \mathbf{z}$ is by construction orthogonal to the space containing $m,$ it follows that ${\mathbb{E}}[\mathbf{z}\varepsilon ]=0,$ so that $\mathbf{z}$ serves as a valid instrument for $\mathbf{x}.$ Identification then follows from the condition that

equation*[equation* omitted — 219 chars of source]

under which the structural parameter satisfies $\bm{\beta} _{0}=\bm{\Sigma} _{{x}}^{-1} \bm{\Sigma} _{xy},$ where $\bm{\Sigma} _{xy}={\mathbb{E}}[\left( \mathbf{x}- \mathbf{P}_{\mathcal{S}}\mathbf{x}\right) y].$

The first condition $m\in {\mathcal{S}}$ specifies the admissible sources of confounding. The second condition $\bm{\Sigma} _{x}>\mathbf{0}$ requires sufficient variation in the structural regressors after those sources have been removed. We develop estimation technology and inference methods for the structural parameters. A key consideration is the choice of ${\mathcal{S}}$, which could be justified by economic, institutional, or design information. Plausible examples include: (a) a known collection of fixed effects or seasonal components; (b) a symmetry or invariance restriction; (c) specified interactions that exclude the structural linear index; (d) a finite-dimensional nuisance space motivated by the empirical design. We explore some examples below. In the absence of such concrete choices, we propose to use the LASSO method applied to a dictionary of basis functions to parsimoniously represent this space.

Our approach is related to several strands of the literature that develop identification strategies derived from instruments internally generated from features of the data generating process. RR2003 shows that heteroskedasticity across regimes can identify structural parameters in simultaneous equations models even in the absence of conventional instruments, see also sp2001. In that framework, changes in the variance of structural shocks generate additional moment conditions that identify the structural coefficients. AL2012 proposes another approach for simultaneous equations based on heteroskedasticity, showing that covariance restrictions between regressors and structural disturbances can generate valid instruments even when external instruments are unavailable.

These methods exploit variation in second moments of the data generating process to produce internal instruments. lewis2025 points to a related literature in macroeconometrics that exploits non-Gaussianity in the structural shocks to obtain identification, such as the independent component analysis (ICA). The literature on endogeneity in general is of course very extensive. There are several well known survey articles by swy2002, anderson2005, sy2005, Imbens2014, ass2019, cr2020, Angrist2023, and lewis2025, and many results have been written into textbooks, such as jmw2001, wooldridge2016, greene2018, and jsmw2018.

The approach developed in this paper follows a similar philosophy but relies on a different source of identification. Instead of exploiting heteroskedasticity or regime variation, we use a projection-based decomposition of the conditional mean of the structural error. This decomposition allows us to construct instrumental variables directly as functions of the regressors. Identification therefore arises from the structure of the conditional expectation rather than from variance shifts or heteroskedasticity--based moment conditions.\footnote{ However, under joint Gaussianity, the endogeneity will be linear and our identification strategy fails as well.}

Our approach is also related to the control function literature discussed, for example, in npv1999, np2003, blundell2004, bck2007, IN2009, newey2013, and JMW2015. These existing methods address endogeneity in semiparametric models by introducing auxiliary instruments or structural restrictions. By contrast, the proposed SP method constructs instruments internally through an orthogonal decomposition of the structural error. This yields a closed--form instrumental variable in each case, which is termed as internally “Constructed Instrumental Variable" (CIV), that depends only on observed regressors and allows identification in a number of classes of linear and nonlinear models as discussed in Appendices (ref) and A.3 below.

In summary, the main distinctions and contributions of this paper are given follows.

\noindenta: \, We propose using a semiparametric projection (SP) method to derive the so--called “CIV" to correctly identify and consistently estimate the unknown parameters and functions involved in a class of correctly identified models under study;

\noindentb: \, We then show that the SP method also offers an alternative way to those methods available from the existing literature about how to deal with endogeneity issues involved in certain classes of weakly identified models;

\noindentc: \, The proposed SP method is invariant to the degree of endogeneity, including a wide range of weak endogeneity, involved in many classes of nonlinear models;

\noindentd: \, We establish the corresponding estimation theory to show that the proposed SP method produces consistent and unbiased estimates as well as asymptotically normal distributions;

\noindente: \, We develop a Hausman--type of nonparametric statistic for testing a full level of strong exogeneity versus a wide range of weak endogeneity with a near--optimal rate;

\noindentf: \, The proposed SP method offers a closed--form expression for an IV function of the data for the correctly identified model under study, and it is easily implementable; and

\noindentg: \, We propose a simple LASSO selection method to evaluate and implement the estimation method by simulated and real data examples.

The organization of this paper is as follows. Sections (ref) and (ref) respectively propose and discuss how to identify and estimate the parameters of interest for a class of correctly identified linear models. Section (ref) then shows that the proposed SP method can also be applied to a class of weakly identified linear models, including such linear models associated with weak endogeneity. Section (ref).1 proposes a LASSO selection method for an optimal set of orthonormal series before it is demonstrated in an analysis of the return to schooling case study in Section (ref).2. Section (ref) concludes and discusses the applicability of the SP method to several classes of extended models.

Appendix (ref) provides some heuristics about the identification and estimation proposed in Sections (ref) and (ref). Appendix (ref) shows that the proposed SP method can be generalized to deal with a wide class of nonlinear models. Appendix A.3 outlines how to extend the proposed SP method to two classes of non-- and semi--parametric regressions, and one class of nonlinear and non--separable binary models. Appendix A.4 briefly discusses other estimation issues including a type of GMM estimation.

Appendix B.1 discusses some implementation issues. Appendix B.2 evaluates the finite--sample properties of the SP method using a number of simulated datasets. Appendix B.3 further investigates the empirical example considered in Section (ref).2 and then adds an extra empirical dataset. Appendix B.4 proposes a jackknife bias--correction method. Appendix C collects the proofs of the main results. Mathematical proofs for Appendix (ref), and additional simulations and extensions, are given in Appendices D-F of an online supplementary document.

For notational consistency, we use $\mathbf{x}$ to stand for a vector regardless of whether it is stochastic or deterministic, and $x$ to stand for either a univariate variable or a given point. When there are no confusions, $x$ can also be used as a generic variable for integration. The usage of several other notation and symbols is standard, such as using $ \mathbf{A}$ as either a matrix or a vector, $\|\mathbf{A}\|^2$ as the standard Euclidean norm, $a$ as a real number, and $[a]\leq a$ as the largest integer part of $a$. There are two modes of convergence involved in the rest of this paper, with “$\rightarrow_{\mathcal{D}}$ and $\rightarrow_P $ representing convergence in distribution and convergence in probability, respectively.

Linear Model Identification

Economic Interpretation and Ex Ante Specification of the Nuisance Space

The identifying content of the model comes from the specification of the nuisance space $\mathcal{S}$. The condition $m(\cdot )\in \mathcal{S}$ should be interpreted as a restriction on the mechanisms through which the endogenous component of the disturbance may depend on the observed state. An ex ante specification of $\mathcal{S}$ should satisfy three principles. First, its defining features should be motivated by economic, institutional, or design information available independently of the outcomes. Second, the restriction should have substantive content: it must exclude at least some functions of the observed state that could otherwise be confused with the structural index $\mathbf{x}^{\top }\bm{\beta }_{0}$. Third, the restriction must leave enough residual variation in $\mathbf{x}$ for ${\mathbf{Q}}_{ \mathcal{S}}=\mathbf{x}-\mathbf{P}_{\mathcal{S}}\mathbf{x}$ to be nonsingular. Ex ante specification does not require every element of the projection operator to be known before seeing the sample. The abstract space $\mathcal{S}$ may be fixed while its inner products are estimated from the observed regressors. For example, the researcher may specify in advance that $\mathcal{S}$ is generated by a set of interactions, calendar functions, or invariant functions, and then use the empirical distribution to orthonormalize those functions or estimate $\mathbf{P}_{\mathcal{S}}\, \mathbf{x}$. Because these operations use only the regressors, they implement the maintained restriction rather than select it according to its implications for the outcome. We next give several suggestive examples.

Example 1. Group and Time Effects. Suppose that observations are indexed by unit $i$ and period $t$, and consider the structural panel model $ y_{it}=\mathbf{x}_{it}^{\top }\bm{\beta}_{0}+\varepsilon _{it}.$ Institutional knowledge may suggest that the conditional mean of the disturbance consists of a unit-specific component and a common time component ${\mathbb{E}}[\varepsilon _{it}\mid x_{it}]=\alpha _{i}+\lambda _{t}.$ The nuisance space is then

equation[equation omitted — 108 chars of source]

This space is specified by the panel structure rather than selected from the outcome data. The residualized regressor is the familiar within-transformed variable, $\mathbf{z}_{it}=\mathbf{x}_{it}-\mathbf{P}_{\mathcal{S}}\,\mathbf{ x}_{it},$ and identification requires nondegenerate variation in $\mathbf{x} _{it}$ after removing the unit and time components: ${\mathbb{E}}[\mathbf{z} _{it}\,\mathbf{z}_{it}^{\top }]>\mathbf{0}$. The economic restriction is not simply that unit and time effects should be included because they improve fit. It is that all systematic dependence between the disturbance and the observed state operates through permanent unit heterogeneity and aggregate time shocks. If, for example, the conditional mean disturbance also contains unit-specific trends that are correlated with $\mathbf{x}_{it}$, the proposed space is misspecified. Such departures can be accommodated through the approximate restriction developed in the paper. A similar type of setting arises in high-frequency demand, production, energy, or financial data, where the institutional environment may imply that predictable omitted shocks follow a known calendar structure described by dummy variables or trigonometric functions, which can similarly be handled in this setting.

Example 2. Symmetry and Invariance Restrictions. In some cases, economic reasoning implies that the confounding component is invariant under a known transformation, whereas the structural component is not. Suppose first that $x$ is scalar and that the selection mechanism depends on the magnitude of $x$, but not on its sign: $m(x)=m(-x).$ Then one may take

equation[equation omitted — 90 chars of source]

The structural function $x\,\beta _{0}$ changes sign under $x\mapsto -x$, whereas every element of $\mathcal{S}_{\mathrm{even}}$ is invariant. Hence, $ \mathrm{span}\{x\}\cap \mathcal{S}_{\mathrm{even}}=\{0\},$ provided the support contains sufficiently rich positive and negative values. The substantive interpretation could be that participation, selection, or measurement error depends on the absolute magnitude of an exposure, while the structural effect depends on its direction. For example, monitoring intensity might respond to the absolute size of a position, whereas the structural return depends on whether the position is long or short. Note that the distribution of $x$ need not itself be symmetric. If it is asymmetric, $P_{\mathcal{S}_{\mathrm{even}}}\,x$ need not be zero, so the residualization is nontrivial and $E[x\,m(x)]$ may be nonzero. Identification nevertheless follows from the separation between the even nuisance functions and the odd structural index.

A multivariate version could be based on permutation invariance. Let $ x=x_{1}-x_{2},$ and suppose that $m(x_{1},x_{2})=m(x_{2},x_{1}).$ The nuisance depends only on the unordered pair of characteristics, while the structural regressor is directional and changes sign when the two components are exchanged. The relevant space is $\mathcal{S}_{\mathrm{sym}}=\left\{ s:s(x_{1},x_{2})=s(x_{2},x_{1})\right\} .$ This structure may arise in matched-pair, bilateral, ranking, or relative-performance settings in which common selection depends symmetrically on the characteristics of the pair, but the parameter of interest multiplies a directional difference.

More generally, if a group of transformations $\mathcal{G}$ acts on the observed state, one may let $\mathcal{S}$ consist of functions invariant under that group: $\mathcal{S}_{\mathcal{G}}=\left\{ s:s(gx)=s(x)\text{ for every} \ g\in \mathcal{G}\right\} .$ Identification is possible when the structural regressors contain components that transform differently from the invariant nuisance functions.

Example 3. Specified Nonlinearities and Interactions. Economic theory may imply that the omitted component operates through known nonlinear features of the observed state but does not contain an unrestricted additive linear component. For example, with two centred regressors, suppose that $m_{0}(x_{1},x_{2})=\gamma _{1}q_{1}(x)+\gamma _{2}q_{2}(x)+\gamma _{3}q_{3}(x),$ where $q_{1}(x)=x_{1}x_{2},$ $ q_{2}(x)=x_{1}^{2}-E[x_{1}^{2}],$ and $q_{3}(x)=x_{2}^{2}-E[x_{2}^{2}].$ The nuisance space is $\mathcal{S}=\mathrm{span}\{q_{1},q_{2},q_{3}\},$ whereas the structural index is $x_{1}\beta _{01}+x_{2}\beta _{02}.$ A possible economic interpretation is that the omitted state reflects complementarities, congestion, dispersion, or adjustment costs that depend on products and squared deviations, while the coefficients of interest measure additive marginal effects. In a production setting, for example, an omitted utilization component might be known to respond to the interaction of capital and labour or to deviations from normal operating levels, rather than to unrestricted linear combinations of the inputs. The basis functions may be orthogonalized with respect to the distribution of $x$. Such orthogonalization can be estimated using the regressors alone and does not change the substantive restriction.

Example 4. Design-Based Strata and Known Assignment Rules. Suppose treatment or exposure varies within strata determined by a known assignment rule. Let $A$ be a predetermined score or collection of baseline characteristics, and suppose the conditional mean disturbance is constant, or follows a specified low-dimensional form, within predefined strata: $m(A)=\sum_{j=1}^{J}\alpha _{j}\,1\{A\in \mathcal{A}_{j}\}.$ The sets $\mathcal{A}_{1},\ldots ,\mathcal{A}_{J}$ might be administrative eligibility groups, sampling strata, geographic zones, cohorts, or categories fixed by the treatment-assignment process. Then $\mathcal{S}= \mathrm{span}\left\{ 1\{A\in \mathcal{A}_{j}\}:j=1,\ldots ,J\right\} .$ The coefficient is identified by variation in $x$ within these predetermined strata: $z=x-{\mathbb{E}}[x\mid \mathcal{A}_{j}]$ for $A\in \mathcal{A}_{j}.$ The identifying assumption is that all systematic selection on unobservables is captured by the known strata.

A related case arises when an administrative assignment algorithm is known to depend on a finite set of basis functions $b_{1}(A),\ldots ,b_{K}(A)$. If economic reasoning implies that the conditional mean of the disturbance depends on the assignment variables through the same features, one may specify $\mathcal{S}=\mathrm{span}\{b_{1}(A),\ldots ,b_{K}(A)\}.$

Example 5. Low-Rank Institutional, Network, or Spatial Components. In some applications, the institutional structure suggests that confounding is generated by a small number of common components. Suppose, for example, that observations belong to known markets or network communities and that $m(x_{i})=\sum_{k=1}^{K}\gamma _{k}f_{k}(x_{i}),$ where the functions $f_{k}$ are fixed before the outcome analysis. They might consist of market indicators, predetermined network eigenvectors, geographic basis functions, or exposure to known aggregate shocks. Then $\mathcal{S}= \mathrm{span}\{f_{1},\ldots ,f_{K}\}.$ Identification comes from the component of the structural regressor that is not explained by these common institutional factors.

Motivations and heuristics

We return to the general structural linear model

equation[equation omitted — 173 chars of source]

Appendix (ref) shows that the proposed SP method is applicable to several classes of nonlinear models, including the general nonparametric regression: $y = g(\mathbf{x}) + \varepsilon$.

Noting that ${\mathbb{E}}[\mathbf{x}\varepsilon]\ne 0$ implies ${\mathbb{E}} [\varepsilon|\mathbf{x}]\ne 0$, we define

equation[equation omitted — 184 chars of source]

where $\bm{\beta}_0$ can be identifiable and consistently estimable as discussed in the rest of this section.

For the nonlinear model ((ref)), we consider the projection process of ${ \varepsilon} = {\mathbb{E}}[{\varepsilon}|\mathbf{x}] + {e}$ in ((ref)) as the proposed the SP method, which has a similar spirit to the so--called “Control Function Approach" discussed in newey1990, npv1999, np2003, JMW2015 and others, but the proposed SP method addresses endogeneity through an internally constructed IV system for such classes of correctly identified models, without assuming the existence and validity of an external IV structure.

To explain the main ideas and for notational simplicity, the main sections of this paper then focus on a semiparametric linear regression of the form:

equation[equation omitted — 78 chars of source]

A diagram is given in Appendix (ref).1 to illustrate model ((ref)) from a geometric point of view. Before we discuss how to identify $\bm{\beta} _0$, we make the following comments.

rem(i) \, Model ((ref)) is quite different from the following partially linear model: \begin{equation} y= \mathbf{x}_u^{\top} \bm{\alpha}_0 + q(\mathbf{x}_v) + \varepsilon, \end{equation} where $\mathbf{x}_u$ and $\mathbf{x}_v$ are observed variables, as discussed in the relevant literature by pmr1988, linton1995, hlg2000, and other more recent developments, for which an identifiability condition: \begin{equation} {\mathbb{E}}\left[(\mathbf{x}_u - {\mathbb{E}}[\mathbf{x}_u|\mathbf{x}_v]) ( \mathbf{x}_u - {\mathbb{E}}[\mathbf{x}_u|\mathbf{x}_v])^{\top}\right]> \mathbf{0} \end{equation} is required for $(\bm{\alpha}_0, q(\cdot))$ to be identifiable. (ii) For model ((ref)) where $\mathbf{x}_u=\mathbf{x}_v = \mathbf{x}$, however, such an identifiability condition in ((ref)) is violated. We therefore show in the rest of Section (ref) that the SP method constructs ${\mathbf{z}}= {\mathbf{x}} - \sum_{j=1}^{\infty} {\mathbb{E}}[{ \mathbf{x}} \, \psi_j(\mathbf{x})] \, \psi_j(\mathbf{x})$ and proposes to replace ((ref)) by requiring ${\mathbb{E}}\left[\mathbf{z} \, \mathbf{z }^{\top}\right]>\mathbf{0}$, where $\{\psi_j(\cdot): j\geq 1\}$ is an orthonormal sequence chosen by the user. (iii) Appendix A.3.1 below discusses that model ((ref)) itself may also be endogenous. In this case, we show that model ((ref)) itself, along with the proposed SP method, is applicable to deal with certain endogeneity issues involved in ((ref)) and some other nonlinear models discussed in Appendix A.3.

To address the endogeneity issue, a simple and commonly used IV method considers combining model ((ref)) with a linear decomposition of $\mathbf{ x}$ of the form:

equation[equation omitted — 147 chars of source]

where $\mathbf{z}_{}$ is assumed to be a valid IV. Under the invertibility of ${\mathbb{E}}\left[\mathbf{z} \, \mathbf{x}^{\top}\right]$, $\bm{\beta}_0$ can be identifiable by $\bm{\beta}_0 = \left({\mathbb{E}}\left[\mathbf{z} \, \mathbf{x}^{\top}\right]\right)^{-1} {\mathbb{E}}[\mathbf{z}_{} \, y]$.

As shown in equations ((ref))--((ref)) below, the proposed SP method consequently enables us to construct an explicit form of $\mathbf{z}$ for it to be the so--called “CIV" as a valid IV.

We now outline the main heuristics about how to construct instrumental variables by the so--called “SP method” to enable that $(\bm{\beta}_0, m(\cdot))$ is identifiable in the rest of this section.

Assume that $m(\mathbf{x})$ can be written as $m(\mathbf{x}) = \mathbf{v}( \mathbf{x})^{\top} \bm{\gamma}$, where $\bm{\gamma}=\left(\gamma_1, \cdots, \gamma_k\right)^{\top}$ is a vector of unknown coefficients, and

equation[equation omitted — 152 chars of source]

is a vector of series functions, where $k$ is a fixed and finite truncation parameter at this stage. The case of varying $k$ will be discussed from Section (ref).2. There is a substantial literature about series estimation, such as gallant1981, andrews1991, newey1997 , ac2003, chen2007, and BCCK2015. A recent book by dg2025 provides an update about nonparametric series estimation.

Letting $\mathbf{v}=\mathbf{v}(\mathbf{x})$ satisfy ${\mathbb{E}}[\mathbf{v} \mathbf{v}^{\top}]>0$, we rewrite model ((ref)) as

equation[equation omitted — 99 chars of source]

where ${\mathbb{E}}[e|\mathbf{x}]=0$. Let $w = y - \mathbf{v}^{\top} \left({ \mathbb{E}}[\mathbf{v} \mathbf{v}^{\top}]\right)^{-1} {\mathbb{E}}[\mathbf{v} y]$. Model ((ref)) can then be written as

equation[equation omitted — 68 chars of source]

where $\mathbf{z} = \mathbf{x} - {\mathbb{E}}[\mathbf{x} \mathbf{v}^{\top}] \left({\mathbb{E}}[\mathbf{v} \mathbf{v}^{\top}]\right)^{-1} \mathbf{v}$ satisfies ${\mathbb{E}}[e|\mathbf{z}] = {\mathbb{E}}[{\mathbb{E}}[e|\mathbf{x }]|\mathbf{z}]=0$, which, along with ${\mathbb{E}}[\mathbf{z} \mathbf{z} ^{\top}]>\mathbf{0}$, ensures that $\mathbf{z}$ can be chosen as an IV. It will be shown in Section (ref).2 that

equation[equation omitted — 135 chars of source]

is identifiable. We therefore have constructed a version of equation ((ref)) of the form:

equation[equation omitted — 344 chars of source]

which implies that $\mathbf{z}$ is the so--called CIV as a valid IV for the linear model: $y = \mathbf{x}^{\top} \bm{\beta}_0 + \varepsilon$.

Assume that $\{(\mathbf{x}_i, y_i)\}$ is a sequence of observed random variables. We introduce $\mathbf{v}_i = \mathbf{v}(\mathbf{x}_i)$, $ w_i = y_i - \mathbf{v}_i^{\top} \left(\sum_{l=1}^n \mathbf{v}_l \, \mathbf{v} _l^{\top}\right)^{-1} \, \sum_{j=1}^n y_j \,\mathbf{v}_j^{\top}$ and $ \mathbf{z}_i = \mathbf{x}_i - \sum_{j=1}^n \mathbf{x}_j \, \mathbf{v} _j^{\top}\, \left(\sum_{l=1}^n \mathbf{v}_l \, \mathbf{v}_l^{\top} \right)^{-1} \mathbf{v}_i$. A sampling version of model ((ref)) is given by:

equation[equation omitted — 92 chars of source]

where $e_i$ is the error term satisfying ${\mathbb{E}}[e_i|\mathbf{x}_i]=0$ and $0<{\mathbb{E}}[e_i^2|\mathbf{x}_i] = \sigma_i^2<\infty$ for $1\leq i\leq n$. Equation ((ref)) implies that $\bm{\beta}_0$ of ((ref)) can be estimated by

equation[equation omitted — 193 chars of source]

where $\mathbf{Z} = \left(\mathbf{z}_1, \cdots, \mathbf{z}_n\right)^{\top}$ and $\mathbf{W} = \left(w_1, \cdots, w_n\right)^{\top}$.

It is pointed out that $\widehat{\bm{\beta}}_{\mathrm{SP}}$ is the same as the following 2sLS estimator:

equation[equation omitted — 201 chars of source]

defined by the conventional two--stage LS method when choosing $\mathbf{Z}$ as the IV function, where $\widetilde{\mathbf{Y}} = \left(\mathbf{I}_n - \mathbf{P}_v\right) \, \mathbf{Y}$ with $\mathbf{Y} = \left(y_1, \cdots, y_n\right)^{\top}$, $\widetilde{\mathbf{X}} = \left(\mathbf{I}_n - \mathbf{P} _v\right) \, \mathbf{X}$ with $\mathbf{X} = \left(\mathbf{x}_1, \cdots, \mathbf{x}_n\right)^{\top}$, $\mathbf{P}_v = \mathbf{V} \, \left(\mathbf{V} ^{\top} \mathbf{V}\right)^{-1} \mathbf{V}^{\top}$ and $\mathbf{V} = \left( \mathbf{v}(\mathbf{x}_1), \cdots, \mathbf{v}(\mathbf{x}_n)\right)^{\top}$.

As noted in ((ref)) above, it is reasonable to require that

equation*[equation* omitted — 426 chars of source]

to ensure that $\bm{\beta}_0$ is correctly identifiable, as discussed rigorously in Section (ref).3 below.

Model identification

To present the main ideas about model identification, we introduce a Hilbert space. Let ${\mathcal{L}}^2=\{g(\mathbf{x}): \ {\mathbb{E}}[g^2( \mathbf{x})]<\infty\}$, where the inner product is given by $\langle g_1( \mathbf{x}), g_2(\mathbf{x})\rangle={\mathbb{E}}[g_1(\mathbf{x})g_2(\mathbf{x })]$ from which the induced norm is $\|g(\mathbf{x})\|=\sqrt{{\mathbb{E}} [g^2(\mathbf{x})]}$ for any given $g_1(\mathbf{x}), g_2(\mathbf{x}), g( \mathbf{x})\in {\mathcal{L}}^2$. Equipped with this $\|\cdot\|$, ${\mathcal{L }}^2$ becomes a Hilbert space.

assumptionSuppose that (i) for any $\lambda\in \mathbb{R}^d$, $ \lambda^\top \mathbf{x}\in {\mathcal{L}}^2$; (ii) ${\mathbb{E}}[ \mathbf{x} \, m(\mathbf{x})]\neq 0$ and ${\mathbb{E}}[m(\mathbf{x})]=0$; (iii) there is a proper closed subspace ${\mathcal{S}}\subsetneqq {\mathcal{L}}^2$ such that $m(\mathbf{x})\in {\mathcal{S}}$ but $\lambda^\top \mathbf{x}\not\in { \mathcal{S}}$ for any $\lambda\ne 0$.

Appendix (ref).2 provides a geometric illustration of Assumption (ref).

rem(a) \, Condition (i) is a minimum requirement for $\mathbf{x}$ in the model, and we also in the sequel need ${\mathbb{E}}[(\lambda^\top \mathbf{x})^2]>0$ for any $\lambda\ne 0$, or equivalently, ${\mathbb{E}}[\mathbf{x}\mathbf{x} ^\top]>0$. Condition (ii) is to confirm the existence of endogeneity in model (ref) since ${\mathbb{E}}[\varepsilon \, \mathbf{x}] = {\mathbb{E} }[ \mathbf{x} \, m(\mathbf{x})]\neq 0$, while ${\mathbb{E}}[m(\mathbf{x})]={ \mathbb{E}}(\varepsilon)=0$. Condition (iii) is crucial that separates the linear form $\lambda^\top \mathbf{x}$ and the unknown $m(\mathbf{x})$ by the proper subspace ${\mathcal{S}}$. (b) \, Condition (iii) requires that there are no more linear combinations of $\mathbf{x}$ left in $m(\mathbf{x})$ to ensure that $\bm{\beta}_0$ involved in ${\mathbb{E}}[y|\mathbf{x}] = \mathbf{x}^{\top} \bm{\beta}_0 + m( \mathbf{x})$ can be correctly identified. Section (ref) shows how to deal with such linearity cases when Condition (iii) fails. (c) \, Section (ref) proposes to choose ${\mathcal{S}}$ by the so--called “LASSO selection method" for practical implementation of the proposed SP method by real data examples, and the choice of ${\mathcal{S}}$ allows that some commonly used functional forms of $m(\mathbf{x})$ can be well expanded in the finite--sample evaluations in Appendices B and E of the online supplementary document.

Given the closed subspace ${\mathcal{S}}$, we have an orthogonal decomposition for ${\mathcal{L}}^2$, that is, ${\mathcal{L}}^2={\mathcal{S}} \oplus{\mathcal{S}}^\bot $ where ${\mathcal{S}}^\bot$ is the orthogonal complement subspace of ${\mathcal{S}}$. This decomposition means that for every element $\xi \in {\mathcal{L}}^2$, it can be uniquely decomposed as $ \xi=\xi_1+\xi_2$ where $\xi_1\in{\mathcal{S}}$, $\xi_2\in {\mathcal{S}}^\bot$ , and they are orthogonal, ${\mathbb{E}}(\xi_1\xi_2)=0$. Accordingly, we have two projection mappings ${\mathcal{P}}_{\mathcal{S}}$ and ${\mathcal{M}} _{\mathcal{S}}$ that map every element in ${\mathcal{L}}^2$ into ${\mathcal{S }}$ and ${\mathcal{S}}^\bot$, respectively, defined by the decomposition, that is, ${\mathcal{P}}_{\mathcal{S}}(\xi)=\xi_1$ and ${\mathcal{M}}_{ \mathcal{S}}(\xi)=\xi_2$.

We establish the following lemma; its proof is given in Appendix C.

lemmaIn addition to Assumption (ref), suppose that $ \{\psi_j(\mathbf{x}), j\ge 1\}$ is an orthonormal basis for ${\mathcal{S}}$, where ${\mathbb{E}}[\psi_j(\mathbf{x})\psi_{\ell}(\mathbf{x})]=\delta_{j\ell} $ and ${\mathbb{E}}[\psi_j(\mathbf{x})]=0$ for all $(j,\ell)$. Then, we have (i) ${\mathcal{P}}_{\mathcal{S}}(\xi)=\sum_{j=1}^\infty {\mathbb{E}}[\xi \psi_j(\mathbf{x})] \psi_j(\mathbf{x})$ for any $\xi \in {\mathcal{L}}^2$; (ii) ${\bm\Sigma}_x\equiv{\mathbb{E}}[\mathbf{xx}^\top]-\sum_{j=1}^\infty { \mathbb{E}}[\mathbf{x} \psi_j(\mathbf{x})] {\mathbb{E}}[\mathbf{x} ^\top\psi_j(\mathbf{x})]>\mathbf{0}$.

Lemma (ref) involves the orthonormality under a probability space. Appendix (ref).3 discusses that all mathematical moments in population can be replaced by their sampling versions, such as replacing ${\mathbb{E}} \left[\psi_j(\mathbf{x}) \, \psi_l(\mathbf{x})\right]$ by $\frac{1}{n} \sum_{i=1}^n \psi_j(\mathbf{x}_i) \, \psi_l(\mathbf{x}_i)$ for all $(j,l)$, in practice, when there is no prior knowledge about the distributional structure of the data under study.

For notational simplicity of our derivations in the rest of this section and the derivations in Appendices A, C and D, we assume that $\{\psi_j(\cdot): j\geq 0\}$ is a sequence of orthonormal functions by definition: ${\mathbb{E} }[\psi_j^2(\mathbf{x}_1)]=1$ and ${\mathbb{E}}[\psi_j(\mathbf{x}_1) \, \psi_l(\mathbf{x}_1)]=0$ for $j\neq l$. Appendices B and E indicate that, moreover, the proposed SP method works well numerically in finite--sample simulations without necessarily requiring orthonormality on $ \{\psi_j(\cdot): j\geq 1\}$.

As the closed subspace of ${\mathcal{L}}^2$, ${\mathcal{S}}$ itself, along with the norm in ${\mathcal{L}}^2$, is a Hilbert space too. After specifying the orthonormal basis, we have the expression of the projection mapping ${ \mathcal{P}}_{\mathcal{S}}$; and therefore the expression of ${\mathcal{M}}_{ \mathcal{S}}=\mathcal{I}-{\mathcal{P}}_{\mathcal{S}}$ where $\mathcal{I}$ is the identity operator. More importantly, the positive definiteness of ${\bm \Sigma}_x$ is the corner stone for the identifiability and estimation of $\bm \beta_0$.

It is noteworthy that the requirement ${\mathbb{E}}[\psi_j(\mathbf{x})]=0$ is to tailor ${\mathcal{S}}$ to cater for ${\mathbb{E}}[m(\mathbf{x})]={ \mathbb{E}}[\varepsilon]=0$. Recalling that ${\mathcal{M}}_{\mathcal{S}}( \mathbf{x})\in {\mathcal{S}}^\bot$, it is orthogonal with any element in ${ \mathcal{S}}$, in particular $m(\mathbf{x})$. Multiplying ${\mathcal{M}}_{ \mathcal{S}}(\mathbf{x})$ on both sides of equation (ref) and taking expectation imply

eqnarray[eqnarray omitted — 320 chars of source]

which yields

equation[equation omitted — 325 chars of source]

which is the expression of $\bm{\beta}_0$ in population.

Note that ${\mathcal{M}}_{\mathcal{S}}(\mathbf{x})$ plays an essential role in the identification of $\bm{\beta}_0$ and satisfies the necessary conditions in the construction of an IV as follows:

eqnarray*[eqnarray* omitted — 431 chars of source]

Lemma (ref) below shows that the expression (ref) is invariant to the choice of the basis in ${\mathcal{S}}$; its proof is given in Appendix C.

lemmaThe expression of $\bm{\beta}_0$ in (ref) is invariant to the choice of the basis in ${\mathcal{S}}$.

We have established a rigorous treatment about the proposed SP method for the identification of $\bm{\beta}_0$ before we discuss about how to estimate $\bm{\beta}_0$ in Section (ref) below.

In the rest of this paper, we consider the case where $d$, the dimensionality of $\mathbf{x}$, is finite and fixed. In Sections (ref)-- (ref), we focus on the case where $d$ is small. Appendix B.1 of the supplemental document discusses how to choose the series function involved for the case of $d\geq 2$ can be large but still fixed in practice.

Estimation and Inference

Estimation method and theory

Let us start to consider estimation and inference issues for model ((ref) ). Suppose that $\{\psi_j(\mathbf{x}), j\ge 1\}$ is an orthonormal basis for ${\mathcal{S}}$, where ${\mathbb{E}}[\psi_j(\mathbf{x})]=0$ and ${\mathbb{E}} [\psi_j(\mathbf{x})\psi_{\ell}(\mathbf{x})]=\delta_{j\ell}$ for all $(j,\ell) $. Hence, we have an orthogonal series expansion for $m(\mathbf{x})$ of the form:

equation[equation omitted — 92 chars of source]

Given a truncation parameter $k$, let $\mathbf{V}_k(\mathbf{x})=(\psi_1( \mathbf{x}), \cdots, \psi_k(\mathbf{x}))^\top$ and $\bm{\gamma}=(\gamma_1, \cdots, \gamma_k)^\top$, and then define $\delta_k(\mathbf{x} )=\sum_{j=k+1}^\infty \psi_j(\mathbf{x}) \, \gamma_j$. Model ((ref)) is written as $y=\mathbf{x}^\top {\bm \beta}_0 +\mathbf{V}_k(\mathbf{x})^\top { \bm \gamma}+\delta_k(\mathbf{x})+e$.

Moreover, given a sample $\{(y_i,\mathbf{x}_i), 1\leq i \leq n\}$, we have $ y_i=\mathbf{x}_i^\top{\bm{\beta}_0}+\mathbf{V}_k(\mathbf{x}_i)^\top {\bm \gamma}+\delta_k(\mathbf{x}_i)+e_i$ for $1\leq i\leq n$. Accordingly, we have a matrix model of the form:

equation[equation omitted — 112 chars of source]

where $\mathbf{y}=(y_1, \cdots, y_n)^\top$, $\mathbf{X}=(\mathbf{x}_1, \cdots, \mathbf{x}_n)^\top$, $\mathbf{V}=(\mathbf{V}_k(\mathbf{x}_1), \cdots, \mathbf{V}_k(\mathbf{x}_n))^\top$, $\mathbf{e}=(e_1, \cdots, e_n)^\top$, and $\bm{\delta}=(\delta_k(\mathbf{x}_1), \cdots, \delta_k( \mathbf{x}_n))^\top$.

Let $\mathbf{P}_v=\mathbf{V}\left(\mathbf{V}^{\top} \mathbf{V}\right)^{-1} \mathbf{V}^{\top}$ and $\mathbf{M}_v= \mathbf{I}_n - \mathbf{P}_v$. We then have

equation*[equation* omitted — 140 chars of source]

Assuming that the matrix involved is invertible, the SP estimator is then given by

align[align omitted — 182 chars of source]

Note that if ${\mathcal{S}}$ is of finite dimension, then ${\bm \delta}=0$ and hence $\widehat{\bm{\beta}}$ is unbiased. For the theoretical generalization, we treat the dimension of ${\mathcal{S}}$ as infinity in the rest of this paper.

In order to establish asymptotic consistency and normality results for $ \widehat{\bm{\beta}}$, we need to introduce the following assumption.

assumptionSuppose that (i) $\{\mathbf{x}_i, i=1, \cdots,n\}$ is a sequence of independent and identically distributed (i.i.d.) observations, and $0<\lim_{n\rightarrow \infty} \frac{1}{n} \sum_{i=1}^n \sigma_i^{2}<\infty$ with $\sigma_i^2={\mathbb{E}}[e_i^2|\mathbf{x}_i]$ (a.s.) for $1\leq i\leq n$; (ii) $\{\psi_j(\cdot)\}$ is chosen such that $ \sup_j{\mathbb{E}}[\psi_j^4(\mathbf{x})]<\infty$; (iii) $m(\cdot)$ is continuously differentiable up to $s\geq d$ order; and (iv) $k$ is a positive integer chosen such that $k^2/n=o(1)$ and $k^{-2s/d} \, n=o(1)$ when $(k, n)\rightarrow (\infty, \infty)$.
remCondition (i) allows for heteroskedasticity. Condition (ii) is easily verifiable for classes of orthonormal series. The smoothness condition imposed on (iii) is also reasonable, which, together with Condition (iv), removes some residue terms in the derivation of asymptotic normality.
theo{Let Assumptions (ref) and (ref) hold.} (i) $\widehat{\bm{\beta}}_{\mathrm{SP}}$ given in (ref) is asymptotically unbiased. (ii) Let $\lim_{n\rightarrow \infty} \frac{1}{n} \sum_{i=1}^n {\mathbb{E}} [e_i^4|\mathbf{x}_i]<\infty$ (a.s.) and ${\mathbb{E}}\left[\left\|\mathbf{x} _i\right\|^4\right]<\infty$. We then have \begin{equation} \sqrt{n} \, \left(\frac{1}{n}\mathbf{X}^\top \mathbf{M}_v \mathbf{\Omega}\mathbf{M}_v\mathbf{X} \right)^{-1/2} \left(\frac{1}{n}\mathbf{X}^\top \mathbf{M}_v\mathbf{X} \right)(\widehat{ \bm{\beta}}_{\mathrm{SP}} -\bm{\beta}_0)\to_{\mathcal{D}} N\left(\mathbf{0}, \mathbf{I}_d\right), \end{equation} as $(k, n)\to (\infty, \infty)$, where $\mathbf{\Omega}=\mathrm{diag} (\sigma_1^2, \cdots, \sigma_n^2)$. (iii) As $(k, n)\to (\infty, \infty)$ and $k^4/n\rightarrow 0$, we have $ \frac{1}{n}\mathbf{X}^\top \mathbf{M}_v\widehat{\mathbf{\Omega }}\mathbf{M}_v\mathbf{X}=\frac{1}{n}\mathbf{X}^\top \mathbf{M} _v\mathbf{\Omega}\mathbf{M}_v\mathbf{X}+o_P(1)$, where $\widehat{\mathbf{ \Omega}}=\mathrm{diag}(\widehat{e}_1^2, \cdots, \widehat{e}_n^2)$, $\widehat{ e}_i=y_i-\mathbf{x}_i^\top\widehat{\bm{\beta}}_{\mathrm{SP}}- \mathbf{V}_k( \mathbf{x}_i)^\top \widehat{\bm \gamma}$ and $\widehat{\bm \gamma}=(\mathbf{V }^\top \mathbf{V})^{-1}\mathbf{V}^\top (\mathbf{y}-\mathbf{X}\widehat{ \bm{\beta}}_{\mathrm{SP}})$.

The proof of Theorem (ref) is given in Appendix C. As shown in part of the proof, we have $\frac{1}{n}\mathbf{X}^\top \mathbf{M}_v \mathbf{\Omega}\mathbf{M}_v\mathbf{X}\rightarrow_P \bm{\Sigma}_x \, \overline{\sigma}_e^{2}$ and $\frac{1}{n}\mathbf{X}^\top \mathbf{M}_v\mathbf{X}\rightarrow_P \bm{\Sigma}_x$ as $(k, n)\to (\infty, \infty)$. Thus, the asymptotic covariance matrix is $\bm{\Sigma}_x^{-1} \, \overline{\sigma}_e^{2}$.

rem(i) Note that the positive definiteness of $\bm{\Sigma}_x$ implied by Lemma (ref) ensures the identifiability of $\bm{\beta}_0$ in ((ref) ). As shown in Lemma (ref) above, meanwhile, $\bm{\Sigma}_x$ is invariant to the choice of $\{\psi_j(\cdot): j\geq 1\}$. It can be seen that $\widehat{\bm{\beta}}$ is asymptotically consistent, and the rate of convergence remains square-root--$n$ even under endogeneity. (ii) Due to $\bm{\Sigma}_x= {\mathbb{E}}[\mathbf{x}\, \mathbf{x}^{\top}] - \sum_{j=1}^{\infty} {\mathbb{E}}[\mathbf{x} \, {\psi}_j(\mathbf{x})] \, { \mathbb{E}}[\mathbf{x}^{\top} {\psi}_j(\mathbf{x})] \leq {\mathbb{E}}[ \mathbf{x}\, \mathbf{x}^{\top}]$, $\widehat{\bm{\beta}}$ is less efficient than the standard OLS method associated when there is no endogeneity. However, $\widehat{\bm{\beta}}$ is the most efficient estimator under Gaussianity on $\{e_i\}$. As shown in Appendix A.4.1 of the supplementary document, moreover, $\widehat{\bm{\beta}}$ is more efficient than $\widehat{ \bm{\beta}}_{\mathrm{LS}}$. (iii) Letting $\widehat{m}(\mathbf{x}) =\mathbf{V}_k(\mathbf{x})^{\top} \widehat{\bm{\gamma}}$ with $\widehat{\bm{\gamma}} = \left(\widehat{\gamma} _1, \cdots, \widehat{\gamma}_k\right)^{\top} = \left(\mathbf{V}^{\top} \mathbf{V}\right)^{-1} \mathbf{V}^{\top} (\mathbf{y} - \mathbf{X} \, \widehat{\bm{\beta}})$, it is pointed out that the availability of $\widehat{ m}(\mathbf{x})$ facilitates the construction of a simple test for checking a full level of exogeneity versus a wide range of weak endogeneity, as discussed in Section (ref).2 below. (v) It is also pointed out that the SP method enables us to obtain $\widehat{ \varepsilon}_i = y_i - \mathbf{x}_i^{\top} \widehat{\bm{\beta}}$ without relying on external instruments, and the availability of $\widehat{ \varepsilon}_i$ may help to reveal possible sources of endogeneity in empirical analysis, as discussed in Appendix B.3.

To start our discussion on the weak endogeneity setting, we propose a simple statistic to test for a full level of exogeneity versus a wide range of weak endogeneity in Section (ref).2.

Testing for weak endogeneity

Consider $y =\mathbf{x}^{\top}\bm{\beta}_0 + m(\mathbf{x}) + e$ with $m( \mathbf{x})={\mathbb{E}}[\varepsilon|\mathbf{x}]$ under the following null hypothesis:

equation[equation omitted — 77 chars of source]

We approximate $m(\mathbf{x})$ by $\sum_{j=1}^k \psi_j(\mathbf{x}) \gamma_{j} \equiv \mathbf{V}_k(\mathbf{x})^{\top} \bm{\gamma}$, where $ \mathbf{V}_k(\cdot) = \left(\psi_1(\cdot), \cdots, \psi_k(\cdot)\right)^{\top}$, and ${\bm\gamma} =\left(\gamma_{1}, \cdots, \gamma_{k}\right)^{\top}$ can be estimated by $\widehat{\bm\gamma}$ in the same way as in Section (ref).1.

In order to test $H_0: \, {\mathbb{P}}(m(\mathbf{x}) =0)=1$, it suffices to test ${\bm\gamma}=0$, or equivalently $(\mathbf{V}^{\top} \mathbf{V}) \, {\bm \gamma}=0$. Since the ${\bm\gamma}$ parameter can be estimated by $\widehat{ \bm\gamma} = (\mathbf{V^{\top} V})^{-1} \mathbf{V}^{\top} (\mathbf{y} - \mathbf{X}\widehat{\bm\beta})$, we introduce $\widetilde{\bm\gamma}_n = ( \mathbf{V^{\top} V}) \, \widehat{\bm\gamma}$ and then define for each given point $x$:

equation[equation omitted — 272 chars of source]

This motives us to propose a simple test statistic (a nonparametric version of the Hausman test proposed in hausman1978) of the form:

eqnarray[eqnarray omitted — 402 chars of source]

by orthonormality: $\int_{\mathcal{X}} \mathbf{V}_k({x}) \mathbf{V}_k^{\top}( {x}) \, f_{\mathbf{x}}(x) d {x}=\mathbf{I}_k$, where $f_{\mathbf{x}}(\cdot)$ denotes the density of $\mathbf{x}$.

In order to avoid dealing with robustness issues regarding the choice of individual $k$ values in practice, we propose using a summarized version of $ L_n(k)$ of the form:

equation[equation omitted — 197 chars of source]

where $k_{\min}$ and $k_{\max}$ are the respective smallest and largest integers satisfying $1\leq k_{\min} < k_{\max}<[\sqrt{n}]$, and $k_{\max} - k_{\min}\rightarrow \infty$ and $\frac{k_{\max}^2}{n}\rightarrow 0$ as $ n\rightarrow\infty$.

theoLet Assumptions (ref) and (ref) hold. Let also $0< \lim_{n\rightarrow \infty} \frac{1}{n} \sum_{i=1}^n {\mathbb{E}}[e_i^4| \mathbf{x}_i]<\infty$ (a.s.) and ${\mathbb{E}}\left[\left\|\mathbf{x} _i\right\|^4\right]<\infty$. We then have under $H_0: \, {\mathbb{P}}(m( \mathbf{x}) =0)=1$: \begin{equation} \frac{T_n}{S_n}\rightarrow_{\mathcal{D}} N(0,1) \ \ \ as $n\rightarrow \infty$, \end{equation} where $S_n^2\equiv 2 \widetilde{\sigma}_e^4 \, \sum_{i=1}^n \sum_{j=1}^n \left(\sum_{k=k_{\mathrm{\min}}}^{k_{\mathrm{\max}}} \mathbf{V}_k^{\top}( \mathbf{x}_i) \mathbf{V}_k(\mathbf{x}_j)\right)^2$, in which $\widetilde{ \sigma}_e^2 = \frac{1}{n} \sum_{i=1}^n \widetilde{e}_i^2$.

Note that $\frac{S_n^2}{\sigma_n^2}\rightarrow_P 1$ with $\sigma_n^2 = \frac{ 2 \,{\sigma}_e^4}{3} \cdot n^2 \, k_{\max}^3\, (1+ o(1))$ as shown in the proof of Theorem (ref) in Appendix C. To show the consistency of our testing statistic under a sequence of local alternatives, we test

equation[equation omitted — 88 chars of source]

where positive sequence $a_n\to 0$ with certain rate while $m(\mathbf{x})\in {\mathcal{L}}^2$ and ${\mathbb{E}}[m^2(\mathbf{x})]>0$.

It is pointed out that such a sequence of local alternatives naturally covers a wide range of weak endogeneity. The following theorem establishes the consistency of the proposed test under $H_1$; its proof is outlined in Appendix C.

theo(i) Let the conditions of Theorem (ref) hold. (ii) Let also $m(\mathbf{x})\in {\mathcal{L}}^2$ with ${\mathbb{E}}[m^2(\mathbf{x})]>0 $. Consider $H_1: \, {\mathbb{P}}(m_n(\mathbf{x})=a_n\, m(\mathbf{x}))=1$, where $a_n\to 0$ and $a_n^2 \, n \, k_{\max}^{-1/2}\to \infty$ as $ n\rightarrow \infty$. We then have under $H_1$, \begin{equation} \frac{\displaystyle1}{\sigma_n} \, T_n\rightarrow_{P} \infty \ \ \ as $n\rightarrow \infty$. \end{equation}

Note that Theorem (ref) shows that the proposed test is capable to detect certain types of weak endogeneity of an order of $a_{n}$, as long as $ a_n$ satisfies $a_{n}^2 \, n \, k_{\max}^{-1/2} \rightarrow \infty$.

Note also that the fast possible rate of $a_n = n^{-\frac{1}{2}} \, k_{\max}^{c_1}$ is near--optimal when $k_{\max} \rightarrow \infty$ at the slowest possible rate, such as $k_{\max} = \left[ c_2 \, \log(\log(n))\right] $, for $c_1>\frac{1}{4}$ and $c_2>0$ to be chosen by the user, although the conventional optimal rate of $a_n = n^{-1/2}$ (see, for example, ss1997) chosen in the parametric setting is not achievable under the proposed SP framework.

In Appendix (ref) below, we show how to extend the proposed SP method to a class of nonlinear models. Appendix A.3 of the supplementary document outlines the main ideas about how to extend the SP method to deal with some other classes of nonlinear models, including two classes of non-- and semi--parametric regression models, and one class of binary models associated with nonlinearity and endogeneity.

Linear Models under Weak Endogeneity

The existing literature pays particular attention on the case where $ (\varepsilon, \mathbf{x})$ in ((ref)) follows a joint Gaussian distribution, which implies that $m(\mathbf{x}) = \left(\mathbf{x} - { \mathbb{E}}[\mathbf{x}]\right)^{\top} \bm{\gamma}_0$ and model ((ref)) becomes $y =\mathbf{x}^{\top} \bm{\beta}_0 + \left(\mathbf{x} - {\mathbb{E}}[ \mathbf{x}]\right)^{\top} \bm{\gamma}_0+ e$ with ${\mathbb{E}}[e|\mathbf{x} ]=0$. Appendix A.4.3 of the main supplementary document discusses that the SP method is applicable to consistently estimate $\bm{\beta}_0$ in certain linearity cases. We now focus on several weak endogeneity settings. When $m( \mathbf{x}) = \left(\mathbf{x} - {\mathbb{E}}[\mathbf{x}]\right)^{\top} \bm{\gamma}_0$, we rewrite model ((ref)): $y = \mathbf{x}^{\top} \bm{\beta}_0 + \varepsilon$ as

equation[equation omitted — 250 chars of source]

which implies that $\left(\bm{\beta}_0 + \bm{\gamma}_0\right) = {\mathbb{E}} ^{-1}\left[\left(\mathbf{x} - {\mathbb{E}}[\mathbf{x}]\right) \left(\mathbf{x } - {\mathbb{E}}[\mathbf{x}]\right)^{\top}\right] \, {\mathbb{E}}\left[\left( \mathbf{x} - {\mathbb{E}}[\mathbf{x}]\right)\, (y - {\mathbb{E}}[y])\right]$ can be correctly identified collectively, rather than $\bm{\beta}_0$ individually.

When there are some additional restrictions imposed on $\bm{\gamma}_0$, such as those weak endogeneity scenarios discussed below, however, $\bm{\beta}_0$ may still be asymptotically identifiable. In such cases where $\bm{\gamma}_0 \equiv \bm{\gamma}_{n0} \rightarrow \mathbf{0}$ as $n\rightarrow \infty$, $ \bm{\beta}_0$ can then be estimated consistently.

We start with model ((ref)) associated with a type of weak endogeneity of the form:

equation[equation omitted — 152 chars of source]

where ${\mathbb{E}}[e_i|\mathbf{x}_i]=0$ for $1\leq i\leq n$, $m_n(\mathbf{x} )={\mathbb{E}}[\varepsilon_i|\mathbf{x}_i = \mathbf{x}] = a_n \, m(\mathbf{x} )$ with ${\mathbb{E}}[m(\mathbf{x})]=0$, ${\mathbb{E}}[\mathbf{x} \, m( \mathbf{x})] \neq 0$, ${\mathbb{E}}[m^2(\mathbf{x})]>0$, $m(\mathbf{x}) \in { \mathcal{S}}$, $a_n\rightarrow 0$ and $a_n \, n\, k^{-\frac{2s}{d} }\rightarrow 0$ as $n\rightarrow \infty$, in which $s$ is the smoothness order of $m(\cdot)$ as in Assumption (ref)(iv). We have as $ n\rightarrow \infty$

equation[equation omitted — 190 chars of source]

Probably because of the projection in ((ref)) and then the definition of weak endogeneity in ((ref)), we are able to show that $\bm{\beta}_0$ can still be consistently estimated.

Since we assume that $m(\mathbf{x}) \in {\mathcal{S}}$, we approximate $m( \mathbf{x})$ in the same way as in Section (ref). In view of model ((ref)), the `endogenous component' represented by $\mathbf{v}^{\top} \bm{\gamma}_n$, with $\bm{\gamma}_n= a_n \, \bm{\gamma}$, has been eliminated. Consequently, equations ((ref)) and ((ref)) show that $ \bm{\beta}_0$ of ((ref)) is correctly identified and therefore it can still be estimated by $\widehat{\bm{\beta}} =\widehat{\bm{\beta}}_{ \mathrm{SP}}$ of ((ref)) consistently.

The establishment and the proof of Theorem (ref) imply Theorem (ref) below.

theoLet the conditions of Theorem (ref)(ii) hold for model ((ref)) with the last part of Assumption (ref)(iv) being weakened to $a_n \, n\, k^{-\frac{2s}{d}}\rightarrow 0$. We then have as $n\rightarrow \infty$ \begin{equation} \sqrt{n} \, \widehat{\bm{\Sigma}}_n(\mathbf{x}, k) \, (\widehat{\bm{\beta}}_{ \mathrm{SP}} - \bm{\beta}_0 ) \rightarrow_{\mathcal{D}} N\left(\mathbf{0}, \, \mathbf{I}_d \right), \end{equation} where $\widehat{\bm{\Sigma}}_n(\mathbf{x}, k) = \widehat{\bm{\Sigma}}_x(k) \, \widehat{\bm{\Sigma}}_n^{-1/2}(k)$, in which $\widehat{\bm{\Sigma}}_x(k) = \frac{1}{n} \sum_{i=1}^n \mathbf{z}_i \, \mathbf{z}_i^{\top}$ and $ \widehat{\bm{\Sigma}}_n(k) = \frac{1}{n} \sum_{i=1}^n \mathbf{z}_i \, \mathbf{z}_i^{\top} \, \widehat{e}_i^2$ with $\widehat{e}_i = w_i - \mathbf{z }_i^{\top} \, \widehat{\bm{\beta}}_{\mathrm{SP}}$, with $(w_i, \mathbf{z}_i)$ being the same as defined in Section 2.1.

Theorem (ref), along with Theorem (ref), shows that the proposed SP method is therefore invariant to the degree of endogeneity under study.

Let us return to the linear weak endogeneity case where $m_n(\mathbf{x}) = \frac{\bm{\mu}^{\top} \left(\mathbf{x} - {\mathbb{E}}[\mathbf{x}]\right)}{ \sqrt{n}}$ with $\|\bm{\mu}\|^2<\infty$. There is a substantial literature as reviewed in a recent survey by ass2019 about treatments of weak endogeneity involving weak IVs.

Let $\bm{\beta}_{n, \bm{\mu}} \equiv: \bm{\beta}_0 + s_n^{-1} \, \sigma_n \, \frac{\bm{\mu}}{\sqrt{n}}$. It follows that the standard OLS estimator becomes

equation[equation omitted — 316 chars of source]

where $m_n(\mathbf{x}) =\frac{\bm{\mu}^{\top} \left(\mathbf{x} - {\mathbb{E}} [\mathbf{x}]\right)}{\sqrt{n}}$, $s_n = \sum_{i=1}^n \mathbf{x}_i \, \mathbf{ x}_i^{\top}$ and $\sigma_n = \sum_{i=1}^n \mathbf{x}_i \, \left(\mathbf{x}_i - {\mathbb{E}}[\mathbf{x}_i]\right)^{\top}$.

We then establish the following theorem; its proof follows trivially from verifying Lindeberg conditions needed for a central limit theorem for an i.i.d. data setting.

theoLet Assumption 3.1(i) hold. If, in addition, ${\mathbb{E}}\left[\left\| \mathbf{x}_i\right\|^4\right]<\infty$ and ${\mathbb{E}}\left[\mathbf{x}_i \, \mathbf{x}_i^{\top}\right]>\mathbf{0}$, then we have \begin{equation} \sqrt{n} \, \widehat{\bm{\Sigma}}_n(\mathbf{x}) \, (\widehat{\bm{\beta}}_{ \mathrm{LS}} - \bm{\beta}_0 ) - {\mathbb{E}}^{-1} [\mathbf{x} \, \mathbf{x} ^{\top} ] \, \mathrm{var}[\mathbf{x}] \, \bm{\mu} \rightarrow_{\mathcal{D}} N\left(\mathbf{0}, \, \mathbf{I}_d \right) \end{equation} as $n\rightarrow \infty$, where $\mathrm{var}[\mathbf{x}] = {\mathbb{E}} [\left(\mathbf{x} - {\mathbb{E}}[\mathbf{x}]\right) \, \left(\mathbf{x} - { \mathbb{E}}[\mathbf{x}]\right)^{\top} ]>\mathbf{0}$ and $\widehat{\bm{\Sigma} }_n(\mathbf{x}) =\widehat{\bm{\Sigma}}_{\mathbf{x}} \, \widehat{\bm{\Sigma}} _e^{-1/2}$, in which $\widehat{\bm{\Sigma}}_{\mathbf{x}} = \frac{1}{n} \sum_{i=1}^n \mathbf{x}_i \, \mathbf{x}_i^{\top}$ and $\widehat{\bm{\Sigma}} _e = \frac{1}{n} \sum_{i=1}^n \mathbf{x}_i \, \mathbf{x}_i^{\top} \, \widetilde{e}_i^2$ with $\widetilde{e}_i = y_i - \mathbf{x}_i^{\top} \widehat{\bm{\beta}}_{\mathrm{LS}}$.

Equation ((ref)) implies that $\widehat{\bm{\beta}}_{\mathrm{LS}} = \bm{\beta}_0 + \widehat{\bm{\Sigma}}_n^{-1} (\mathbf{x})\, {\mathbb{E}}^{-1} [\mathbf{x} \, \mathbf{x}^{\top} ] \, \mathrm{var}[\mathbf{x}] \, \frac{ \bm{\mu}}{\sqrt{n}} + O_P\left(\frac{\widehat{\bm{\Sigma}}_n^{-1}(\mathbf{x}) }{\sqrt{n}}\right) \rightarrow \bm{\beta}_0$, indicating that $\widehat{ \bm{\beta}}_{\mathrm{LS}}$ is a consistent estimator of $\bm{\beta}_0$.

However, $\bm{\mu}$ itself cannot be identified correctly. This becomes clearer in the special case where $\bm{\beta}_0= \mathbf{0}$, $ m^{\ast}(\cdot) =0$ and $y = \left(\frac{\bm{\mu}}{\sqrt{n}}\right)^{\top} \left(\mathbf{x} - {\mathbb{E}}[\mathbf{x}]\right) + e$. In this case, we have as $n\rightarrow \infty$

equation[equation omitted — 282 chars of source]

which shows that $\widehat{\bm{\mu}} = \sqrt{n} \, \left(\sum_{i=1}^n \left( \mathbf{x}_i - \overline{\mathbf{x}}_n\right) \, \left(\mathbf{x}_i - \overline{\mathbf{x}}_n\right)^{\top}\right)^{-1} \, \sum_{i=1}^n \left( \mathbf{x}_i - \overline{\mathbf{x}}_n\right) \, y_i$ is not consistent to $ \bm{\mu}$, although $\frac{\widehat{\bm{\mu}}}{\sqrt{n}}$ is consistent to $ \frac{\bm{\mu}}{\sqrt{n}}$, where $\overline{\mathbf{x}}_n = \frac{1}{n} \sum_{i=1}^n \mathbf{x}_i$.

This is similar in spirit to those discussed in the weak IV literature, see equation (2.5) of ss1997, for example. We demonstrate in Appendix B.4 of the supplementary document that one may eliminate $\bm{\mu}$ by a simple jackknife method.

Furthermore, we show that the SP method offers a simple way to deal with a mixture of strong and weak endogeneity of the form:

equation*[equation* omitted — 156 chars of source]

which implies

equation*[equation* omitted — 191 chars of source]

where ${\mathbb{E}}[m^{\ast}(\mathbf{x})]=0$, ${\mathbb{E}}|m^{\ast}(\mathbf{ x})|^2>0$ and $m^{\ast}(\mathbf{x}) \in {\mathcal{S}}$. Our approach allows either ${\mathbb{E}}[\mathbf{x} \, m^{\ast}(\mathbf{x})] = 0$ or ${\mathbb{E} }[\mathbf{x} \, m^{\ast}(\mathbf{x})] \neq 0$.

We use the same $\{\psi_j(\mathbf{x}): j\geq 1\}$ as an orthonormal series to approximate $m^{\ast}(\mathbf{x})$ by $m_{k}^{\ast}(\mathbf{x}) \equiv:\sum_{j=1}^{k} \psi_j(\mathbf{x}) \, \gamma_{j}^{\ast}$, where $ \bm{\gamma}^{\ast}_k = \left(\gamma_{1}^{\ast}, \cdots, \gamma_{k}^{\ast}\right)^{\top}$ is an array of coefficients.

In the same spirit as in Section (ref).1, we estimate $\bm{\beta}_{n, \bm{\mu}}= \bm{\beta}_0 + s_n^{-1} \, \sigma_n \, \frac{\bm{\mu}}{\sqrt{n}}$ by $\widehat{\bm{\beta}}_{\mathrm{SP}}$ of ((ref)). Obviously, $\bm{\beta} _0$ is only asymptotically identifiable. We finally establish the following theorem; its proof follows trivially from that of Theorem (ref).

theoLet Assumptions (ref) and (ref) hold with $m(\cdot)$ being replaced by $m^{\ast}(\cdot)$. We have as $n\rightarrow \infty$ \begin{equation} \sqrt{n} \, \widehat{\bm{\Sigma}}_n(\mathbf{x}, k) \, (\widehat{\bm{\beta}}_{ \mathrm{SP}} - \bm{\beta}_0 ) - {\mathbb{E}}^{-1} [\mathbf{x} \, \mathbf{x} ^{\top} ] \, \mathrm{var}[\mathbf{x}] \, \bm{\mu} \rightarrow_{\mathcal{D}} N\left(\mathbf{0}, \, \mathbf{I}_d\right). \end{equation}

Equation ((ref)) implies $\widehat{\bm{\beta}}_{\mathrm{SP}} = \bm{\beta}_0 + \widehat{\bm{\Sigma}}_n^{-1}(\mathbf{x}, k) \, {\mathbb{E}} ^{-1} [\mathbf{x} \, \mathbf{x}^{\top} ] \, \mathrm{var}[\mathbf{x}] \, \frac{\bm{\mu}}{\sqrt{n}} + O_P\left(\frac{\widehat{\bm{\Sigma}}_n^{-1}( \mathbf{x}, k)}{\sqrt{n}}\right) \rightarrow \bm{\beta}_0$, which indicates that $\widehat{\bm{\beta}}_{\mathrm{SP}}$ is still a consistent estimator of $\bm{\beta}_0$.

Section (ref) proposes a simple LASSO selection method to realise the proposed SP method in practice before we employ it for a real dataset analysis in Section (ref) and then in Appendix B.2 of the supplementary document for the finite--sample evaluation of simulated datasets and another real dataset.

LASSO Selection and Empirical Analysis

Construction of $\mathcal{S}$

Constructing the subspace ${\mathcal{S}}$ involved in Assumption 2.1 is crucial in our theory that affects both the asymptotic consistency and unbiasedness of the estimator proposed. We focus on the linear model discussed in Sections 2 and 3, and a similar discussion for the nonlinear setting in Appendix A.2 follows accordingly.

As can be seen, since the limit covariance matrix of $\widehat{\bm \beta}_{ \mathrm{SP}}$ of (3.2) is proportional to ${\bm \Sigma}_x^{-1}$ where ${\bm \Sigma}_x={\mathbb{E}}[\mathbf{x}\mathbf{x}^\top]-\sum_{j=1}^\infty {\mathbb{ E}}[\mathbf{x}\psi_j(\mathbf{x})]{\mathbb{E}}[\psi_j(\mathbf{x})\mathbf{x} ^\top]$, in which the smaller the number of basis functions we use, the more efficient $\widehat{\bm \beta}_{\mathrm{SP}}$ can be. Therefore, it is not unreasonable that we choose such $\mathcal{S}$ that forms the smallest subspace containing $m(\mathbf{x})$. We then propose a data driven approach to selecting the basis functions that spans the subspace $\mathcal{S}$.

Note that, as a closed subspace of the Hilbert space $\mathcal{L}^2$, $ \mathcal{S}$ is a Hilbert space too. Since $\mathcal{L}^2$ is separable, it possesses an orthonormal sequence as its basis (see, Page 169 of dudley2003, for example); some subsequence of this basis will be the basis of $\mathcal{S}$ by construction. Precisely, let $\{\phi_j(\mathbf{x}), j\ge 0\}$ be the basis of $\mathcal{L}^2$ and the subsequence $\{\phi_{s_j}(\mathbf{x}), j=1, 2, \cdots\}$ be the basis of $\mathcal{S}$ that we denote as $\{\psi_j(\mathbf{x }), j=1, 2, \cdots\}$ in previous sections, that is, $\psi_j(\mathbf{x} )=\phi_{s_j}(\mathbf{x})$ for $j=1,2 \cdots$. The choice of $\{\phi_j( \mathbf{x}), j\ge 0\}$ should avoid orthogonal polynomial sequence because the endogeneity implies ${\mathbb{E}}[\mathbf{x}m(\mathbf{x})]\ne 0$; otherwise, this condition cannot be fulfilled. Consequently, with the help of Lasso approach and the regression equation, along with Assumption (ref), we are able to select $\{\psi_j(\mathbf{x}), j=1, 2, \cdots\}$ from its parent sequence $\{\phi_j(\mathbf{x}), j\ge 0\}$, and therefore, $ \mathcal{S}=$span$\{\psi_j(\mathbf{x}), j=1, 2, \cdots\}$.

{To elaborate this data driven method,} without loss of generality, we consider the function space ${\mathcal{L}}^2$ where the support of $\mathbf{x }$ is $[0,\pi]$. Accordingly, we introduce an orthonormal basis in ${ \mathcal{L}}^2$ of the form: $\{\phi_0(\mathbf{x}) =1, \phi_j(\mathbf{x})= \sqrt{2/\pi}\cos( j \mathbf{x}) \text{ for } j\ge 1\}$, which has been employed for the numerical exercises conducted in Section (ref) and Appendix B.2. In Appendix E of the online supplementary document, we also consider some other basis functions to check the robustness of the SP method. In the rest of this subsection, our goal becomes to find the basis of the subspace $\mathcal{S}$ from the basis $\{\phi_j(\mathbf{x}), j\ge 0\}$ in ${\mathcal{L}}^2$.

{Since $m(\mathbf{x})\in \mathcal{S}\subsetneqq{{\mathcal{L}}^2}$,} we may expand $m(\mathbf{x})$ in terms of $\{\phi_j(\mathbf{x}), j\ge 0\}$ by $m( \mathbf{x})=\sum_{j=1}^\infty \phi_j(\mathbf{x}) \, \xi_j$, where we use $ \{\xi_j: j\geq 1\}$ as the coefficients {in order to distinguish the $ \{\gamma_j: j\geq 1\}$ in the expansion of $m(\mathbf{x})$ in terms of $ \{\psi_j(\mathbf{x}), j\ge 1\}$. Here, we abandon the constant term $\phi_0( \mathbf{x})$ due to ${\mathbb{E}}[m(\mathbf{x})]=0$. More importantly, many $ \gamma_j$ are either zero or statistically insignificant. This is because $m( \mathbf{x})$ belongs to the subspace $\mathcal{S}$, and by Parseval's identity in Hilbert space, ${\mathbb{E}}[m(\mathbf{x})^2]=\sum_{j=1}^\infty \xi_j^2<\infty$ implying the attenuation of $\xi_j$.}

{Meanwhile, from the sieve estimation point of view, we could not estimate an infinite dimensional function; chen2007 (see, Page 5561 for example). Instead, we use a “sieve”, $\cdots \subset {\mathcal{L}}^2_k\subset {\mathcal{ L}}^2_{k+1}\subset \cdots \subset {\mathcal{L}}^2$ where ${\mathcal{L}}^2_k=$ span$\{\phi_j(\mathbf{x}), j=1, \cdots, k\}$ that approach ${\mathcal{L}}^2$ ; correspondingly, ${\mathcal{S}}_{(k)}=$span $\{\phi_{s_j}(\mathbf{x}), s_j\le k\}$ form a sieve for ${\mathcal{S}}$. Once we identify each element in ${\mathcal{S}}_{(k)}$ with high probability for each $k$, we obtain $ \widehat{\mathcal{S}}_{(k)}$ by the data driven method.}

{To proceed,} given a sample $\{(\mathbf{x}_i, y_i), i=1,\cdots,n\}$, we can rewrite the model as follows:

equation[equation omitted — 140 chars of source]

where $k$ is the truncation parameter, and $\delta_{k}(\mathbf{x} _i)=\sum_{j=k+1}^\infty \phi_j(\mathbf{x}_i) \, \xi_j$. The parameters $ \{\xi_j\}$ need to meet certain sparsity conditions that will be specified soon.

Based on (ref), we propose the following LASSO selection method.

equation[equation omitted — 292 chars of source]

where $\widehat{\bm\xi}^*=(\widehat{\xi}_1^*,\ldots, \widehat{\xi}_k^*)^\top$ , and $k$ and $\lambda$ are the user--chosen truncation and tuning parameters, respectively.

equation[equation omitted — 308 chars of source]

where $\{\zeta_j\}$ is a set of pre-determined weights such as $1/\widehat{ \xi}_j^*$ with $\widehat{\xi}_j^*$ being the $j^{th}$ element of $\widehat{ \bm\xi}^*$. \ The LASSO estimate $\widehat{\bm\xi}^\dag=(\widehat{ \xi} _{1}^\dag,\ldots, \widehat{ \xi}_{k}^\dag)^\top$ enables us to obtain

eqnarray[eqnarray omitted — 141 chars of source]

for which we show that $\Pr(\widehat{\mathcal{S}}_{(k)}=\mathcal{S} _{(k)})\rightarrow 1$ as $(n, k) \rightarrow (\infty, \infty)$ in the following Theorem (ref).

assumption(i) Let $\mathscr{C}=\{\ell \mid \xi_\ell\ne 0, \, \ell \le k\}$ with cardinality $\sharp \mathscr{C} =k_0$, and $\bar{\mathscr{C}}=\{\ell \mid \xi_\ell = 0, \, \ell \le k\}$ for any given sufficiently large $k$. (ii) There exists a $J> 2$ such that ${\mathbb{E}}\left[|e_i|^J\right]<\infty $, $\frac{k}{[n^{J-1}(\log k)^{J/2}]}\to 0$, and $\frac{k_0 \sqrt{\log k}}{ \sqrt{n}}\to 0$.

Assumption (ref)(i) is a typical sparse restriction and infers that $\mathscr{C}\cap \bar{\mathscr{C}}=\emptyset$ and $\sharp \bar{ \mathscr{C}}=k-k_0$. Assumption (ref)(ii) regulates the truncation parameter and the sample size further. {In some special cases, e.g. $m(\cdot) $ is even function or odd function, $k_0\le k/2$, this condition is fulfilled.}

We establish the following asymptotic consistency; its proof is given in Appendix C.

theoLet the conditions of Theorem 3.1 and Assumption (ref) hold, and let $\lambda\asymp \frac{\sqrt{\log k}}{\sqrt{n}} $. We then obtain as $(k, n)\to (\infty, \infty)$ (i) { $\|\operatorname*{\normalfont\textrm{vec}}(\boldsymbol{\beta}_0 - \widehat{\boldsymbol{\beta}}^*, \boldsymbol{\xi}-\widehat{\boldsymbol{\xi}}^* ) \| =O_P(\frac{\sqrt{k_0\log k }}{\sqrt{n}})$; (ii) \ $\|\operatorname*{\normalfont\textrm{vec}}(\boldsymbol{\beta}_0 - \widehat{\boldsymbol{ \beta}}^*, \boldsymbol{\xi}-\widehat{\boldsymbol{\xi}}^* ) \|_1 =O_P(\frac{ k_0\sqrt{\log k}}{\sqrt{n}})$}. Suppose that { $\min_{\ell \in \mathscr{C}} |\xi_{\ell}| \gg \sqrt{ (k_0\log k )/n} (1+ \max_{\ell\in \mathscr{C}} \zeta_{\ell} )$, and $ \min_{\ell \in \bar{\mathscr{C}}} \zeta_{\ell}\gg k_0 \left(1+ \max_{\ell\in \mathscr{C}} \zeta_{\ell}\right)$}, additionally. (iii) Then { $\Pr( \operatorname*{\normalfont\textrm{sgn}}(\widehat{\pmb{\xi}}_k^{\dag}) =\operatorname*{\normalfont\textrm{sgn}}(\pmb{\xi}))\to 1$. (iv) $\sqrt{n}( \widehat{\bm \beta}^\dag -{\bm \beta}_0)\to_D N(\mathbf{0}, \bm{\Sigma} _x^{-1} \, \overline{\sigma}_e^{2})$, where $\bm{\Sigma}_x^{-1} \, \overline{ \sigma}_e^{2}$} is the same as defined below Theorem (ref).

The additional conditions in Theorem (ref) can be easily fulfilled. It is worth mentioning that $\Pr(\operatorname*{\normalfont\textrm{sgn}}(\widehat{\pmb{\xi}} _k^{\dag}) =\operatorname*{\normalfont\textrm{sgn}}(\pmb{\xi})) \to 1$ implies that $\Pr(\widehat{\mathcal{S}} _{(k)}=\mathcal{S}_{(k)})\rightarrow 1$ as $(k, n)\to (\infty, \infty)$. Consequently, we have chosen the subset $\mathcal{S}_{(k)}$ asymptotically.

The return to schooling case study

\noindentExample 5.1: In this example, we examine the returns to schooling (see, card2001, for example) using data from the 1979 National Longitudinal Survey of Youth\footnote{ The data is collected from \url{https://davidcard.berkeley.edu/data_sets.html}.}. After removing those individuals having missing values, we have $n=2639$ observations left in the dataset. We specifically consider the following variables:

{\

table[table omitted — 435 chars of source]

}

In Table (ref), $y$ is the $\log$(wage) (i.e., the response), and $x$ is the endogenous regressor (i.e., the number of years of child's education received). We follow Chapter 15 of wooldridge2016 as well as agmr2023 and cff2025 to consider the following set of IVs: $z_1$ stands for the number of years of dad's education; $z_2$ equals to 1 if the individual grew up near an accredited four-year college, 0 otherwise; $z_3$ equals to 1 if the individual grew up near an accredited two-year college, 0 otherwise; $z_4 $ equals to 1 if the individual grew up near an accredited four-year public college, 0 otherwise; $z_5$ equals to 1 if the individual grew up near an accredited four-year private college, 0 otherwise.

We consider the following linear model:

equation[equation omitted — 75 chars of source]

and our focus is to infer $\beta_0$. To be consistent with the simulation designs considered in Appendix B, we rescale $x_i$'s as follows:

equation[equation omitted — 107 chars of source]

where $\underline{x} =\min_{i}x_i$ and $\overline{x}=\max_i x_i$. Similarly, we normalize $z_{1i}$'s. This is in the same spirit of the data transformation process as in cff2025. Thus, $\widetilde{x}_i\in [0,\pi]$. After the data transformation (ref), the model becomes

equation[equation omitted — 107 chars of source]

where $\widetilde{\beta}_0 = \frac{1}{\pi}(\overline{x} - \underline{x} )\beta_0$, $\widetilde{\alpha}_0 =\frac{1}{\pi} \underline{x}\beta_0 + \alpha_0$, and $\varepsilon_i =\widetilde{m}(\widetilde{x}_i)+e_i$ with $ \widetilde{m}(\widetilde{x}_i) = m(\widetilde{x}_i \cdot \frac{\overline{x} - \underline{x}}{\pi}+ \underline{x}) =m(x_i)$. By doing so, we maintain the property of $\varepsilon_i$ in order to compare with the 2SLS method widely adopted in the literature. For the purpose of comparison, we consider the following methods:

(i) SP method; \, (ii) LS method: (a) regressing $y$ on $(1, \widetilde{x})$ ; (b) regressing $y$ on $(1, \widetilde{x}, z_{1}, z_{1}^2)$; \, (iii) 2sLS method for the model (ref) under two scenarios: (a) Just identified: using $z_j$ as IV only for $j=1,2,\ldots, 5$ respectively; (b) Over identified: using two IVs $z_{1}$ and $z_{j}$, where $j=2,3,4,5$; (c) Over identified: using $z_{1}$ and $\phi_2(\cdot),\phi_3(\cdot),\phi_8(\cdot)$ together as IVs, where $\phi_2(\cdot)$, $\phi_3(\cdot)$, $\phi_8(\cdot)$ are selected by the SP method; and (d) Over identified: using $\phi_2(\cdot)$, $ \phi_3(\cdot)$, $\phi_8(\cdot)$ as IVs.

In what follows, we consider the case where the level of the dad's education is included in the following model:

equation[equation omitted — 188 chars of source]

where $x$ and $z_1$ denote the number of years of the education level received for the child and dad, respectively. We estimate the model using the SP method as well. We then select $\mathcal{S}$ from $\{\phi_0( \widetilde{x}),\phi_1(\widetilde{x}),\ldots \}\otimes \{\phi_0(z),\phi_1(z),\ldots \}$. Using the proposed LASSO method, we select $\phi_1(\widetilde{x})\phi_6(z)$, $\phi_2(\widetilde{x})$, and $\phi_3( \widetilde{x})\phi_1(z)$. The selection of the basis functions is somewhat consistent with those in the main text, as we also identify $\phi_2( \widetilde{x})$ in the LASSO procedure. The estimated $m(x,z)$ is $\widehat{m }(x,z)= 0.0680\phi_1(\widetilde{x})\phi_6(z)-0.0678 \phi_2(\widetilde{x}) -0.1232 \phi_3(\widetilde{x})\phi_1(z)$. We refer to this model as SP* in what follows, and still focus on reporting $\widetilde{\beta}_0$.

For each method, we report the estimated value of $\widetilde{\beta}_0$ and it 95% confidence interval (referred to CI in Table (ref)). Moreover, we report the in-sample root mean squared errors of each method: $ \text{RMSE} =\sqrt{\frac{1}{n}\sum_{i=1}^n (y_i-\widehat{y}_i)^2}$.

{

table[table omitted — 1,468 chars of source]

}

We summarize the results in Table (ref). All methods confirm a positive relationship between wage and education, and the SP method offers better model fitting in terms of RMSE. The LASSO method selects $ \{\phi_2(x), \phi_3(x), \phi_8(x)\}$. It is noteworthy that Table (ref) provides the estimates of both $\widetilde{\beta}_0$ and $\beta_0$ . After accounting for the relationship $\widetilde{\beta}_0 = \frac{1}{\pi}( \overline{x} - \underline{x})\beta_0$, the estimates of $\beta_0$ are reported in the last column. The result from the SP method shows that an additional year eduction suggests around 7.1% increase in wage on average. The finding from the SP method is very close to 2sLS method when the IVs are selected as $z_1$ only, as $(z_1, z_3)$, or as $(z_1, z_5)$.

Meanwhile, Table (ref) also shows that the proposed SP* model accommodates the case where $\mathbf{x}$ includes externally observed variables, such as the level of the dad's education. Among the SP, SP* and LS(b) estimates, SP* performs better than SP and LS(b), while LS(b) is a commonly used one in such empirical analysis. By using different IVs, obviously, the outcomes of 2sLS vary from case to case, which does create some uncertainty for practical interpretation. By contrast, the SP and SP* estimation offers robust estimates regardless of the choice of the series functions and truncation parameters.

Appendix B.3.2 examines possible source of endogeneity involved in Example 5.1. In Appendix B.3.3, additionally, we discuss possible endogeneity for a stock return dataset and offer a simple and robust way to deal with the case where the structural model error term may be endogenously correlated with the market portfolio variable under study.

Conclusions and Discussion

This paper has proposed the SP method to deal with certain types of endogeneity involved in linear and nonlinear models. The corresponding estimation and testing theory has been established, and evaluated by both simulated and real datasets. Our finite--sample results established in Section (ref).2, Appendix B of the supplementary document and Appendix E of the online appendix have shown that the proposed SP method works well numerically with the support of the sound large--sample theory established in Sections (ref)--(ref).1.

Appendix (ref) below discusses about how to extend the proposed SP method to the identification, estimation and testing for a class of nonlinear regression models. As alluded in the introduction and then outlined in Appendix A.3 of the supplementary document, the proposed SP method can also be extended in a unified way to a number of non-- and semi--parametric regression models associated endogeneity issues.

Since such models have their own identifiability issues, existing developments in the relevant literature, and there are estimation and testing properties to be established for each of them, we wish to write up details for each of them separately.

{

appendix

Linear and Nonlinear Models

Linear regression

A diagrammatic illustration about model ((ref))

In model ((ref)): $y={\bf x}^\top{\bm\beta}_0+\varepsilon$, the regressor ${\bf x}$ has an effect on $y$ in two ways. One is a direct effect that adds ${\bf x}^\top {\bm \beta}_0$ to $y$; and the other is an indirect effect that affects $y$ through $\varepsilon$, that is, when ${\bf x}$ changes, the error term $\varepsilon$ changes as well due to correlation with ${\bf x}$, as illustrated below.

center[center omitted — 389 chars of source]

More precisely, the projection $m({\bf x})={\mathbb E}(\varepsilon|{\bf x})$ splits the contribution of $\varepsilon$ to $y$ into two parts, $\varepsilon=m({\bf x})+e$, where the first contribution $m({\bf x})$ stems from ${\bf x}$ while the second contribution $e$ is due to all factors orthogonal with ${\bf x}$. The effect of $m({\bf x})$ on $y$ is of all factors implicitly included in $\varepsilon$ that are correlated with ${\bf x}$. When ${\bf x}\sim {\bf x}+\Delta {\bf x}$, we have

equation[equation omitted — 264 chars of source]

For the endogenous model, because of these two--fold effect of ${\bf x}$ on $y$, OLS could not identify the direct effect of ${\bf x}$. Our method extracts the “energy” of ${\bf x}$ in $\varepsilon$, exposing the effect $m({\bf x})$ in the equation, so that under certain conditions we can identify the direct and indirect effects of ${\bf x}$ on $y$.

A geometric illustration about Assumption (ref)

We have defined the Hilbert space ${\mathcal L}^2=\{g({\bf x}):\; {\mathbb E}[g^2({\bf x})]<\infty\}$, in which ${\mathbb E}[\xi_1\xi_2]$ is an inner product for $\xi_1, \xi_2\in {\mathcal L}^2$; in Assumption (ref), we stipulate ${\mathcal S}$ as the closed proper subspace of ${\mathcal L}^2$. Accordingly, we have direct sum decomposition ${\mathcal L}^2={\mathcal S}\oplus {\mathcal S}^\bot$ where ${\mathcal S}^\bot$ is the orthogonal complement subspace of ${\mathcal S}$ that is the set of all elements from ${\mathcal L}^2$ orthogonal with each element in ${\mathcal S}$. Notice that both ${\mathcal S}$ and ${\mathcal S}^\bot$ are Hilbert space too, with the norm defined in ${\mathcal L}^2$.

At element level, for any $\xi\in{\mathcal L}^2$, one can uniquely decompose $\xi=\xi_1+\xi_2$ with $\xi_1\in{\mathcal S}$ and $\xi_2\in{\mathcal S}^\bot$. Accordingly, two projection mappings can be defined as ${\mathcal P_S}(\xi)=\xi_1$ and ${\mathcal M_S}(\xi)=\xi_2$. Certainly, ${\mathcal M_S}\equiv {\mathcal I-P_S}$ where ${\mathcal I}$ is the identity operator.

In Assumption (ref), the condition $\lambda^\top {\bf x}\not \in {\mathcal S}$ for any $\lambda\neq 0$ implies that ${\mathcal M_S}(\lambda^\top{\bf x})$ does not degenerate, as illustrated in Figure (ref).

{ \tikzset{global scale/.style={ scale=#1, every node/.append style={scale=#1} } } {

figure[figure omitted — 607 chars of source]

} }

Consequently, as shown in Lemma 2.1 above, we have the following identifiability condition:

equation[equation omitted — 248 chars of source]

ensures that $\bm{\beta}_0 = {\bm\Sigma}_x^{-1} \, \bm{\Sigma}_{xy}$ is identifiable, where $\bm{\Sigma}_{xy} = {\mathbb E}[\mathbf{x} \, y] - \sum_{j=1}^\infty {\mathbb E}[{\bf x} \, \psi_j({\bf x})] \, {\mathbb E}[y \, \psi_j({\bf x})]$. Furthermore, Lemma 2.2 shows that $\bm{\beta}_0$ is invariant to the choice of the basis $\{\psi_j(\cdot): j\geq 1\}$.

Probability density function of $\mathbf{x}$ unknown

Since the space ${\mathcal L}^2$ is defined as the set of all $g({\bf x})$ such that ${\mathbb E}[g^2({\bf x})]<\infty$, when the density of $\mathbf{x}$ is unknown, it is not easy to determine which function of ${\bf x}$ belongs to the space. In this case, we may relax the restriction ${\mathbb E}[g^2({\bf x})]<\infty$.

Let $\pi(x)$ be a density satisfying $c\pi(x)\le f_{\bf x}(x)\le C\pi(x)$ for some constants $C>c>0$, where $f_{\bf x}(x)$ is the density of ${\bf x}$. Define ${\mathcal L}^2_{\pi}=\{h({\bf x}): \int h^2(x)\pi(x)dx<\infty\}$. It is clear that $g({\bf x})\in {\mathcal L}^2$ iff $g({\bf x})\in {\mathcal L}^2_{\pi}$. Let ${\mathcal S}_{\pi}$ be a proper closed subspace of ${\mathcal L}^2_{\pi}$, and ${\mathcal S}_{\pi}={\rm span}\{\xi_j({\bf x}), j\ge 1\}$ where $\{\xi_j(x)\}$ is an orthonormal sequence w.r.t. $\pi(x)$, that is, $\int \xi_i(x)\xi_j(x)\pi(x)dx=\delta_{ij}$.

Because of Proposition 2.1 in BCCK2015, for any vector $U_k({\bf x})=(\xi_1({\bf x}), \cdots, \xi_k({\bf x}))$, the eigenvalues of the matrix ${\mathbb E}[U_k({\bf x})U_k({\bf x})^\top]$ are bounded above and away from zero, uniformly over $k$. This implies that any subset of $\{\xi_j({\bf x}), j\ge 1\} $ is linearly independent, so we use Gram-Schmidt procedure to orthogonalize $\{\xi_j({\bf x}), j\ge 1\} $ to be $\{{\psi}_j({\bf x}), j\ge 1\}$, such that (1) ${{\mathbb E}}[{\psi}_i({\bf x}){\psi}_j({\bf x})]=\delta_{ij}$, and ${\mathbb E}[{\psi}_j({\bf x})]=0$ for all $i, j$; (2) ${\rm span}\{{\psi}_j({\bf x}), j\ge 1\}={\rm span}\{{\xi}_j({\bf x}), j\ge 1\}$.

Therefore, $\{{\psi}_j({\bf x}), j\ge 1\}$ is an orthonormal sequence under inner product $\langle\xi, \eta\rangle={\mathbb E}(\xi \eta)$, and ${\mathcal S}={\rm span}\{{\psi}_j({\bf x}), j\ge 1\}$ is a proper subspace of ${\mathcal L}^2$; moreover, $m({\bf x})\in {\mathcal S}$ iff $m({\bf x})\in {\mathcal S}_\pi$.

We shall show these assertions using induction. First, let ${\psi}_1({\bf x})=\frac{\displaystyle1}{\displaystyle{\rm s.d.}(\xi_1({\bf x})) }\{\xi_1({\bf x})-{\mathbb E}[\xi_1({\bf x})]\}.$ Clearly, the assertion holds for $k=1$, that is, ${\rm span}\{\psi_1({\bf x})\}={\rm span}\{{\xi}_1({\bf x})\}$ and ${\mathbb E}[{\psi}_1({\bf x})]=0$.

Suppose the assertion holds for $k-1$, i.e. $\{{\psi}_j({\bf x}), j=1,\cdots, k-1\}$ is a set of orthonormal sequence, ${\rm span}\{\psi_j({\bf x}), j=1,\cdots, k-1\}={\rm span}\{{\xi}_j({\bf x}), j=1,\cdots, k-1\}$, and ${\mathbb E}[{\psi}_j({\bf x})]=0$, $\forall \, j\le k-1$. Define $h_k({\bf x})=\xi_k({\bf x})-{\mathbb E}[\xi_k({\bf x})]-\sum_{j=1}^{k-1}{\mathbb E}[\xi_k({\bf x}) \psi_j({\bf x})] \, {\psi}_j({\bf x}).$

$h_k({\bf x})$ is orthogonal with each of $\{{\psi}_j({\bf x}), j=1,\cdots, k-1\}$, ${\mathbb E}[h_k({\bf x})]=0$, and ${\mathbb E}[h_k^2({\bf x})]={\rm Var}[\xi_k({\bf x})]-\sum_{j=1}^{k-1}\{{\mathbb E}[\xi_k({\bf x}) \, {\psi}_j({\bf x})]\}^2>0;$ Otherwise $\xi_k({\bf x})$ is a combination of $\{{\psi}_j({\bf x}), j=1,\cdots, k-1\}$, a contradiction. Let ${\psi}_k({\bf x})=\frac{\displaystyle1}{\displaystyle{\rm s.d.}(h_k({\bf x}))} \, h_k({\bf x}).$ Hence, the assertion holds for $k$. By the method of induction, the assertion holds for any given $k$ and therefore for the sequence.

Noting that ${\rm span}\{\psi_j({\bf x}), j=1,\cdots\}={\rm span}\{{\xi}_j({\bf x}), j=1,\cdots\}$, $m({\bf x})$ can be expanded as an orthogonal infinite series in terms of $\{{\psi}_j({\bf x}), j\geq 1\}$ iff it can be expanded as an orthogonal infinite series in terms of $\{{\xi}_j({\bf x}), j\geq 1\}$. So, $m({\bf x})\in {\mathcal S}$ iff $m({\bf x})\in {\mathcal S}_\pi$. It is worth mentioning that in the Gram-Schmidt procedure, all mathematical moments in population can be replaced by their sampling versions, such as replacing ${\mathbb E}\left[\psi_j(\mathbf{x}) \, \psi_l(\mathbf{x})\right]$ by $\frac{1}{n} \sum_{i=1}^n \psi_j(\mathbf{x}_i) \, \psi_l(\mathbf{x}_i)$ for all $(j,l)$, in practice, when there is no prior knowledge about the distributional structure of the data under study.

Nonlinear regression

Estimation of nonlinear regression

Consider a general nonlinear regression model of the form:

equation[equation omitted — 75 chars of source]

where the functional form of $g(\cdot, \bm{\theta}_0)$ is assumed to be parametrically known, but is indexed by ${\bm \theta}_0$ as a vector of unknown parameters, and $(\varepsilon, \mathbf{x})$ is the same as defined in Section (ref), in which the $p$--dimensional ${\bm \theta}_0$ is an interior point of a compact set $\Theta\subset \mathbb{R}^p$.

The relevant literature (see, amemiya1977, amemiya1985, BT1985, newey1990, for example) discusses endogeneity issues for model ((ref)). Before we discuss about how to identify and then estimate $\bm{\theta}_0$, we provide the following examples.

{\bf Example A.1}: Consider the following log(wage) and education model:

equation[equation omitted — 176 chars of source]

where $y$ denotes the log(wage), $x_{1}$, $x_{2}$ and $x_2$ denote the education, experience, and experience$^2$, respectively, and $\varepsilon$ is uncorrelated with $x_2$ and $x_2^2$, but is correlated with $x_1$ (see, for example, wooldridge2016), and $\bm{\theta}_0 = \left(\theta_{00}, \theta_{01}, \theta_{02}, \theta_{03}\right)^{\top}$ is a vector of unknown parameters.

Appendix E shows that the proposed SP estimation method works well numerically when $m(\mathbf{x})$ contains second--order polynomial terms of $\mathbf{x}$.

{\bf Example A.2}: Consider the following additive case where

equation[equation omitted — 160 chars of source]

where $\bm{\theta}_0 = \left(\bm{\alpha}_{01}^{\top}, \cdots, \bm{\alpha}_{0q}^{\top}; \gamma_{01}, \cdots, \gamma_{0q}\right)^{\top}$ is a vector of unknown parameters of interest, and each $g_j(\cdot)$ is a known function commonly used in empirical applications.

Model ((ref)) covers a class of important models often used in the neural network literature for the case where $g_j(\cdot) = \sigma(\cdot)$ is chosen as the so--called activation function (see, for example, BK2019). While permitting possible endogeneity, Example E3 of Appendix E shows that the SP method works well numerically for a special case of model ((ref)).

We now extend our SP method to identify and estimate $\bm{\theta}_0$ in model ((ref)). Recall $m(\mathbf{x}) = {\mathbb E}[\varepsilon|\mathbf{x}]$ and then rewrite model ((ref)) as follows:

equation[equation omitted — 174 chars of source]

Note that the introduction of model ((ref)) not only covers the linear model in ((ref)), but also facilitates the discussion of a wide range of nonlinear and non--separable econometric models as outlined in Appendix A.3 of the supplementary document.

As in model ((ref)), we rewrite model ((ref)) as follows:

equation[equation omitted — 131 chars of source]

where $r(\mathbf{x}) = \sum_{j=k+1}^{\infty} \psi_j(\mathbf{x}) \, \gamma_j$, and the true version: $(\bm{\theta}_0, \bm{\gamma}_0)$ is chosen such that to minimize

equation[equation omitted — 224 chars of source]

Letting $\frac{\partial S(\bm{\theta}, \bm{\gamma})}{\partial \bm{\theta}}|_{(\bm{\theta} = \bm{\theta}_0, \bm{\gamma} = \bm{\gamma}_0)}=0$ and $\frac{\partial S(\bm{\theta}, \bm{\gamma})}{\partial \bm{\gamma}}|_{(\bm{\theta} = \bm{\theta}_0, \bm{\gamma} = \bm{\gamma}_0)}=0$, we then have

eqnarray[eqnarray omitted — 620 chars of source]

which imply

eqnarray[eqnarray omitted — 498 chars of source]

where $\mathbf{g}_1(\mathbf{x}, \bm{\theta}_0) =\frac{\partial g(\mathbf{x}, \bm{\theta})}{\partial \bm{\theta}}|_{\bm{\theta} = \bm{\theta}_0}$.

Recall ${\mathcal L}^2$, ${\mathcal S}$, ${\mathcal P}_{\mathcal S}$ and ${\mathcal M}_{\mathcal S}$ as defined in the same way as in Section 2.2 above. Recall also that the operator ${\mathcal P}_{\mathcal S}$ projects any element of ${\mathcal L}^2$ into ${\mathcal S}$, while the operator ${\mathcal M}_{\mathcal S}$ projects any element of ${\mathcal L}^2$ into ${\mathcal S}^\bot$. In particular, we have ${\mathcal M}_{\mathcal S}(\varepsilon)=\varepsilon-{\mathcal M}_{\mathcal P}(\varepsilon)=e$.

If we use the projection mapping operator, equation ((ref)) reduces to

equation*[equation* omitted — 144 chars of source]

which is the first order condition of $\min_{{\bm \theta}}{\mathbb E}[{\mathcal M_S}(y-\mathbf{g}(\mathbf{x}, \bm{\theta}))]^2$ at $\bm{\theta} = \bm{\theta}_0$. It follows from model ((ref)), ${\mathcal M_S}(y-\mathbf{g}(\mathbf{x}, \bm{\theta}_0))=e$, which is the corresponding nonlinear parametric model without endogeneity. Therefore, the conventional NLS applies. In a similar way to the proof of Lemma (ref) above, it can be shown that the identifiability of $\bm{\theta}_0$ by ((ref)) is invariant to the choice of $\{\psi_j(\cdot): j\geq 1\}$.

In the case of $g(\mathbf{x}, \bm{\theta}_0) = \mathbf{x}^{\top} \bm{\beta}_0$ in Section (ref), Equation ((ref)) reduces to Lemma (ref)(ii) by requiring

equation[equation omitted — 217 chars of source]

to ensure that $\bm{\beta}_0 = {\bm\Sigma}_x^{-1} \, \bm{\Sigma}_{xy}$ is identifiable, where $\bm{\Sigma}_{xy} = {\mathbb E}[\mathbf{x} \, y] - \sum_{j=1}^\infty {\mathbb E}[{\bf x} \, \psi_j({\bf x})] \, {\mathbb E}[y \, \psi_j({\bf x})]$.

We now introduce the following assumption.

assumptionLet ${\mathcal S}$={\rm span}$\{\psi_1(\mathbf{x}), \psi_2(\mathbf{x}), \cdots\}$ where $\{\psi_j(\cdot),\, j\ge 1\}$ is an orthonormal sequence with ${\mathbb E}[\psi_j({\bf x})]=0$ and $\sup_j{\mathbb E}[\psi_j^4({\bf x})]<\infty$. Let ${\mathcal S}_0=\{g(\mathbf{x}, {\bm\theta}), {\bm\theta}\in \Theta\}$. Suppose that (i) $m(\mathbf{x})\in {\mathcal S}$, (ii) ${\mathcal S}_0\cap {\mathcal S}=\emptyset$ or $\{0\}$, and $({\mathcal S}_0-g(\mathbf{x}, {\bm \theta}_0))\cap {\mathcal S}=\{0\}$ where $0$ stands for the null function; and (iii) ${\mathbb E}\left[\mathbf{g}_1(\mathbf{x}, {\bm \theta}_0)\, m(\mathbf{x})\right]\ne 0$.
remNote that ${\mathcal S}_0-g(\mathbf{x}, {\bm \theta}_0)$ means the set of all elements of ${\mathcal S}_0$ minus $g(\mathbf{x}, {\bm \theta}_0)$. Condition (i) confines the function $m(\mathbf{x})$ in ${\mathcal S}$. Due to this, we may be able to identify ${\bm \theta}_0$ in the endogenous model. Note also that, as a set, $\mathcal{S}_0$ does not necessarily contain the null function, while, as a subspace, $\mathcal{S}$ does. In Condition (ii), ${\mathcal S}_0\cap {\mathcal S}=\emptyset$ or $\{0\}$ means that, at most $\mathcal{S}_0$ and $\mathcal{S}$ have a common function, i.e. the null function; the set ${\mathcal S}_0-g(\mathbf{x}, {\bm \theta}_0)$ shifts all elements of ${\mathcal S}_0$ by $g(\mathbf{x}, {\bm \theta}_0)$, and we require the intersection $({\mathcal S}_0- g(\mathbf{x}, {\bm \theta}_0))\cap {\mathcal S}$ only contains the null function that facilitates the establishment of the consistency of the estimator defined later. In addition, Conditions (ii) and (iii) together imply that $g(\mathbf{x}, {\bm \theta}_0)\not \in {\mathcal S}$ and $g(\mathbf{x}, {\bm \theta}_0)\not \in {\mathcal S}^\bot$. On the one hand, this maintains the endogeneity in model (ref), and on the other hand, it enables us to identify ${\bm \theta}_0$. Consequently, Assumption (ref) ensures that we are able to identify ${\bm \theta}_0$ and estimate it by the NLS method. The rest of the discussion of Assumption (ref) is similar to that of Assumption (ref).

Therefore, operating ${\mathcal M_S}$ on both sides of model (ref) yields ${\mathcal M_S}(y)={\mathcal M_S}(g({\bf x},{\bm \theta}_0))+e$. This motivates the following observation:

eqnarray[eqnarray omitted — 416 chars of source]

where the equality holds if and only if ${\mathbb E}\{[{\mathcal M_S}(g({\bf x},{\bm \theta}_0)-g({\bf x},{\bm \theta}))]^2\}=0$. Thus, under Equation ((ref)) and Assumption (ref), ${\bm \theta}_0$ is the unique minimum point of ${\mathbb E}\{[{\mathcal M_S}(y-g({\bf x},{\bm \theta}))]^2\}$, so that ${\bm \theta}_0$ is uniquely identifiable and NLS is applicable to estimate it.

Given $\{(y_i, {\bf x}_i), i=1, \cdots,n\}$, a sampling version of model (ref) is as follows:

equation[equation omitted — 105 chars of source]

Note that, under Assumption (ref), we have $m(\mathbf{x})=\sum_{j=1}^\infty \psi_j(\mathbf{x}) \, \gamma_j$ with $\gamma_j={\mathbb E}[m(\mathbf{x}) \, \psi_j(\mathbf{x})]$, and for a given truncation parameter $k>1$, define the partial sum $m_k(\mathbf{x})=\mathbf{V}_k(\mathbf{x})^\top {\bm\gamma}$, where $\mathbf{V}_k(\mathbf{x})=(\psi_1(\mathbf{x}), \cdots, \psi_k(\mathbf{x}))^\top$ and ${\bm\gamma}=(\gamma_1, \cdots, \gamma_k)^\top$.

We also define $\delta_k(\mathbf{x})=\sum_{j=k+1}^\infty \gamma_j \psi_j(\mathbf{x})$ for us to rewrite ((ref)) as

equation[equation omitted — 135 chars of source]

which can be written in matrix form:

equation[equation omitted — 109 chars of source]

where $\mathbf{G}({\bm \theta}_0)=(g({\bf x}_1, {\bm \theta}_0), \cdots, g({\bf x}_n, {\bm \theta}_0))^\top$.

Defining ${\bf P}_v={\bf V}({\bf V}^\top {\bf V})^{-1}{\bf V}^\top$ and ${\bf M}_v={\bf I}_n-{\bf P}_v$, we then have from (ref) that ${\bf M}_v({\bf y}-\mathbf{G}({\bm \theta}_0))={\bf M}_v({\bm\delta}+{\bf e})$. Letting $L_n( {\bm \theta})=\frac{\displaystyle1}{\Dn}\|{\bf M}_v({\bf y}- \mathbf{G}({\bm \theta}))\|^2$, the estimator $\widehat{\bm\theta}$ is then defined by

equation[equation omitted — 119 chars of source]
theo[Consistency] Suppose $\{(y_i, {\bf x}_i), i=1, \cdots,n\}$ is an i.i.d. sequence, ${\mathbb E}[e_1^2]=\sigma_e^2<\infty$. In addition to Assumption (ref), suppose that for any ${\bm \theta}\in \Theta$, ${\bm \theta}\ne{\bm \theta}_0$, $\lambda(\{x: \; g(x, {\bm \theta})\ne g(x, {\bm \theta}_0)\})>0$ where $\lambda$ is Lebesgue measure; $k^2=o(n)$ as $n\to\infty$. Then, $\widehat{\bm \theta}\to_P{\bm \theta}_0$ as $(k, n)\rightarrow (\infty, \infty)$.

The additional condition $\lambda(\{x: \; g(x, {\bm \theta})\ne g(x, {\bm \theta}_0)\})>0$ is necessary for the asymptotic consistency of the NLS estimator, as it helps identify ${\bm \theta}$ from ${\bm \theta}_0$ when $g(x, {\bm \theta})\ne g(x, {\bm \theta}_0)$ for a set of points, $x$'s, whose measure is greater than zero.

Let $\mathbf{S}_n({\bm \theta})= \frac{\displaystyle\partial}{\displaystyle\partial {\bm \theta}}L_n( {\bm \theta})=-\frac{\displaystyle2}{\Dn}\frac{\displaystyle\partial}{\displaystyle\partial {\bm \theta}} \mathbf{G}({\bm \theta})^\top {\bf M}_v({\bf y}- \mathbf{G}({\bm \theta}))$ and

align*[align* omitted — 564 chars of source]

be the score function and Hessian matrix of $L_n( {\bm \theta})$, respectively. Before we establish Theorem (ref) below, we introduce the following assumption.

assumptionLet $g(\mathbf{x}, {\bm\theta})$ be differentiable w.r.t. ${\bm\theta}$ up to second order such that (i) All elements of $\frac{\displaystyle\partial}{\displaystyle\partial {\bm \theta}} g(\mathbf{x}, {\bm \theta}_0)$ are not in $\mathcal{S}$; (ii) All elements of $\frac{\displaystyle\partial^2}{\displaystyle\partial {\bm \theta} \partial {\bm \theta}^\top}g(\mathbf{x}, {\bm \theta}_0)$ are in ${\mathcal L}^2$; (iii) The Hessian matrix $\mathbf{H}_n({\bm \theta})$ of $L_n({\bm \theta})$ are such that for some sequence $\epsilon_n\to 0$ as $n\to\infty$, $\sup_{\|{\bm \theta}-{\bm \theta}_0\|<\epsilon_n}\| \mathbf{H}_n({\bm \theta}) - \mathbf{H}_n({\bm \theta}_0)\|=o_P(1)$; (iv) Assumption 3.1(iii)(iv) remains satisfied.
rem(a) Conditions (i) and (ii) are commonly used in NLS estimation {while the condition (i) excludes the derivative $\frac{\displaystyle\partial}{\displaystyle\partial {\bm \theta}} g(\mathbf{x}, {\bm \theta}_0)$ from $\mathcal{S}$, likewise $g(\mathbf{x}, {\bm \theta}_0)$}. The condition imposed on (iii) removes some residue terms in the derivation of asymptotic normality. As shown in Theorem 23.3 of oliver2017book, if $\widehat{\bm \theta}\to_P{\bm \theta}_0$, then there exists a sequence $\epsilon_n\to 0$ as $n\to\infty$ such that $\|\widehat{\bm \theta}-{\bm \theta}_0\|<\epsilon_n$ with probability tending to one. Hence, under the consistency of $\widehat{\bm \theta}$, we are able to focus on a shrinking neighbourhood of ${\bm \theta}_0$ in the establishment of an asymptotic normality. In particular, if $\frac{\displaystyle\partial^2}{\displaystyle\partial {\bm \theta} \partial {\bm \theta}^\top}g({\bf x}, {\bm \theta}_0)$ is Lipschitz, the condition (iii) is fulfilled automatically (see, Assumption 3.4 in dgl2023, for example). (b) Note that $\bm{\Sigma}_g = {\mathbb E}\left[\mathbf{g}_1(\mathbf{x}, \bm{\theta}_0) \, \mathbf{g}_1^{\top}(\mathbf{x}, \bm{\theta}_0)\right] - \sum_{j=1}^{\infty} {\mathbb E}\left[\mathbf{g}_1(\mathbf{x}, \bm{\theta}_0) \, \psi_j(\mathbf{x})\right] \, {\mathbb E}\left[\mathbf{g}_1^{\top}(\mathbf{x}, \bm{\theta}_0)\, \psi_j(\mathbf{x})\right]>{\bf 0}$ follows from Assumptions (ref) and (ref), where $\mathbf{g}_1(\mathbf{x}, \bm{\theta}_0)$ is assumed to satisfy the following condition: $\sum_{j=1}^{\infty} \|{\mathbb E}\left[\mathbf{g}_1(\mathbf{x}, \bm{\theta}_0) \, \psi_j(\mathbf{x})\right]\|^2<\infty$.
theo[Normality] Suppose $\{(y_i, {\bf x}_i), i=1, \cdots,n\}$ is an i.i.d. sequence with ${\mathbb E}[e_1^4]<\infty$ and ${\mathbb E}\left[\left\|\frac{\displaystyle\partial}{\displaystyle\partial {\bm\theta}}g({\bf x}_1, {\bm\theta}_0)\right\|^4\right]<\infty$. Under Assumptions (ref) and (ref), we then have as $(k, n)\rightarrow (\infty, \infty)$, \begin{align} \sqrt{n} \, \mathbf{S}_n(\bm \theta_0)\to_{\mathcal D}N(\mathbf{0}, \sigma_e^2{\bm\Sigma}_g) \ \ \ and \ \ \ \mathbf{H}_n(\bm \theta_0) &\to_P{\bm\Sigma}_g, \end{align} where $\bm{\Sigma}_g = {\mathbb E}\left[\mathbf{g}_1(\mathbf{x}, \bm{\theta}_0) \, \mathbf{g}_1^{\top}(\mathbf{x}, \bm{\theta}_0)\right] - \sum_{j=1}^{\infty} {\mathbb E}\left[\mathbf{g}_1(\mathbf{x}, \bm{\theta}_0) \, \psi_j(\mathbf{x})\right] \, {\mathbb E}\left[\mathbf{g}_1^{\top}(\mathbf{x}, \bm{\theta}_0)\, \psi_j(\mathbf{x})\right]>0$. Consequently, \begin{equation} \sqrt{n} \, (\widehat{\bm \theta}-{\bm \theta}_0)\to_{\mathcal D} N(\mathbf{0}, \sigma_e^2 \, {\bm\Sigma}_g^{-1}) \ \ \ as $n\rightarrow \infty$. \end{equation}

The unknown quantities involved in Theorem (ref) can be consistently estimated respectively by their sampling versions. It can be shown that $\bm{\Sigma}_g$ is invariant to the choice of $\{\psi_j(\cdot): j\geq 1\}$ in an analogous way to Lemma 2.2. We also have the following remarks.

rem(i) \, Our discussion about model ((ref)) covers an extended version of the form: \begin{equation} y = g(\mathbf{x}, \bm{\theta}_0) + \varepsilon = g(\mathbf{x}_{\rm exo}, \mathbf{x}_{\rm end}; \bm{\theta}_0) + \varepsilon, \ \ with \ ${\mathbb E}[\varepsilon|\mathbf{x}_{\rm end}]\neq 0$ \ but \ ${\mathbb E}[\varepsilon|\mathbf{x}_{\rm exo}]=0$, \end{equation} where $\mathbf{x}_{\rm end}$ and $\mathbf{x}_{\rm exo}$ are the endogenous and exogenous components of $\mathbf{x}$, respectively. (ii) \, Note that when we specify ${\mathbb E}[\varepsilon|\mathbf{x}] =m(\mathbf{x}, \bm{\gamma}_0)$ parametrically, model ((ref)) then becomes a parametrically nonlinear model of the form: \begin{equation} y = g(\mathbf{x}, \bm{\theta}_0) + \varepsilon = g(\mathbf{x}, \bm{\theta}_0) + m(\mathbf{x}, \bm{\gamma}_0) + e, \end{equation} where $(\bm{\theta}_0, \bm{\gamma}_0)$ is identifiable and estimable when \begin{eqnarray} \bm{\Sigma}_{gm} \equiv && {\mathbb E}\left[\mathbf{g}_1(\mathbf{x}, \bm{\theta}_0) \, \mathbf{g}_1^{\top}(\mathbf{x}, \bm{\theta}_0)\right] \nonumber\\ && - {\mathbb E}\left[\mathbf{g}_1(\mathbf{x}, \bm{\theta}_0) \, \mathbf{m}_1^{\top}(\mathbf{x}, \bm{\gamma}_0)\right] \bm{\Sigma}_{mm}^{-1} {\mathbb E}\left[\mathbf{m}_1(\mathbf{x}, \bm{\gamma}_0)\, \mathbf{g}_1^{\top}(\mathbf{x}, \bm{\theta}_0)\right]>{\bf 0}, \end{eqnarray} in which $\bm{\Sigma}_{mm} = {\mathbb E}\left[\mathbf{m}_1(\mathbf{x}, \bm{\gamma}_0)\, \mathbf{m}_1^{\top}(\mathbf{x}, \bm{\gamma}_0)\right]>{\bf 0}$, $\mathbf{g}_1(\mathbf{x}, \bm{\theta}_0) =\frac{\partial g(\mathbf{x}, \bm{\theta})}{\partial \bm{\theta}}|_{\bm{\theta} = \bm{\theta}_0}$ and $\mathbf{m}_1(\mathbf{x}, \bm{\gamma}_0) =\frac{\partial m(\mathbf{x}, \bm{\gamma})}{\partial \bm{\gamma}}|_{\bm{\gamma} = \bm{\gamma}_0}$. It can be shown that equation ((ref)) is required for the identifiability of $(\bm{\theta}_0, \bm{\gamma}_0)$. In a similar way to that of ${\bm\Sigma}_g$ involved in Theorem (ref), the corresponding covariance matrix becomes \begin{equation} \mathbf{\Omega}_{gm} \equiv: \begin{pmatrix} {\mathbb E}\left[\mathbf{g}_1(\mathbf{x}, \bm{\theta}_0) \, \mathbf{g}_1^{\top}(\mathbf{x}, \bm{\theta}_0)\right] & {\mathbb E}\left[\mathbf{g}_1(\mathbf{x}, \bm{\theta}_0) \, \mathbf{m}_1^{\top}(\mathbf{x}, \bm{\gamma}_0)\right] \\ {\mathbb E}\left[\mathbf{m}_1(\mathbf{x}, \bm{\gamma}_0) \, \mathbf{g}_1^{\top}(\mathbf{x}, \bm{\theta}_0)\right] & {\mathbb E}\left[\mathbf{m}_1(\mathbf{x}, \bm{\gamma}_0) \, \mathbf{m}_1^{\top}(\mathbf{x}, \bm{\gamma}_0)\right] \end{pmatrix}^\top, \end{equation} which is invertible under Condition ((ref)). (iii) \, Specifically, we now show that the proposed SP approach is applicable to a class of nonlinear models of the form: \begin{equation} y = g(\mathbf{x}, \bm{\theta}_0) + \varepsilon = g(\mathbf{x}, \bm{\theta}_0) + m(\mathbf{x}) + e = g(\mathbf{x}, \bm{\theta}_0) + \mathbf{x}^{\top} \bm{\gamma}_0 + e, \end{equation} where $\bm{\theta}_0$ is identifiable when $\bm{\Sigma}_{gn} \equiv : {\mathbb E}\left[\mathbf{x} \, \mathbf{x}^{\top}\right] - {\mathbb E}\left[\mathbf{x} \, \mathbf{g}_1^{\top}(\mathbf{x}, \bm{\theta}_0)\right] \bm{\Sigma}_{gg}^{-1} {\mathbb E}\left[\mathbf{g}_1(\mathbf{x}, \bm{\theta}_0)\, \mathbf{x}^{\top}\right]$ is positive definite, in which $\mathbf{g}_1(\mathbf{x}, \bm{\theta}_0) =\frac{\partial g(\mathbf{x}, \bm{\theta})}{\partial \theta}|_{\bm{\theta} = \bm{\theta}_0}$ and $\bm{\Sigma}_{gg} = {\mathbb E}\left[\mathbf{g}_1(\mathbf{x}, \bm{\theta}_0)\, \mathbf{g}_1^{\top}(\mathbf{x}, \bm{\theta}_0)\right]$ is invertible, and $\bm{\gamma}_0$ is an unknown parameter. Therefore, model ((ref)) shows that the proposed SP method covers the nonlinear regression case where $(\varepsilon, \mathbf{x})$ follows a joint Gaussian distribution, and $g(\mathbf{x}, \bm{\theta}_0)$ is nonlinear in $\mathbf{x}$. As discussed above, model ((ref)) covers that $(\varepsilon, \mathbf{x})$ follows a joint Gaussian distribution. By slightly modifying the assumptions and these proofs of Theorems (ref) and (ref) for model ((ref)), a corresponding estimation theory can be established accordingly for model ((ref)).

In view of the above discussion in Section (ref).2 on weak endogeneity, we also propose a test for a full level of exogeneity versus a wide range of nonlinear weak endogeneity.

Testing for nonlinear weak endogeneity

Consider $y =g(\mathbf{x}, \bm{\theta}_0) + \varepsilon = g(\mathbf{x}, \bm{\theta}_0) + m_n(\mathbf{x}) + e$ under the following null hypothesis:

equation[equation omitted — 70 chars of source]

with $m_n(\mathbf{x})={\mathbb E}[\varepsilon|\mathbf{x}]$ satisfies ${\mathbb E}[m_n^2({\bf x})]\rightarrow 0$ as $n\rightarrow \infty$.

To test $H_0: \, {\mathbb P}(m_n({\bf x}) =0)=1$, in a similar fashion to that discussed in Section (ref), we can establish an asymptotic normality of $T_n$ that is defined in the same way as in ((ref)) in Section (ref) with $\widetilde{e}_i = y_i - g({\bf x}_i, \widehat{\bm\theta})$.

theoLet Assumptions (ref) and (ref) hold. Let also ${\mathbb E}\left[\left\|\frac{\displaystyle\partial}{\displaystyle\partial {\bm\theta}}g({\bf x}_1, {\bm\theta}_0)\right\|^4\right]<\infty$ and ${\mathbb E}[e_1^4]<\infty$. We then have under $H_0: \, {\mathbb P}(m_n({\bf x}) =0)=1$: \begin{equation} \frac{\displaystyle1}{\sigma_n} \, T_n\rightarrow_{\mathcal D} N(0,1) \ \ \ as $n\rightarrow \infty$, \end{equation} where $\sigma_n^2\equiv 2 \widehat{\sigma}_e^4 \, \sum_{i=1}^n \sum_{j=1}^n {\mathbb E}\left[\left(\sum_{k=k_{\rm \min}}^{k_{\rm \max}} \mathbf{V}_k^{\top}({\bf x}_i) \mathbf{V}_k({\bf x}_j)\right)^2\right]$, in which $\widehat{\sigma}_e^2 = \frac{1}{n} \sum_{i=1}^n \widetilde{e}_i^2$.

Equation ((ref)) shows that $T_n$ is asymptotically normally distributed under $H_0$. The proof of Theorem (ref) is given in Appendix C, where it can also be shown $\sigma_n^2= \frac{2 \,\widehat{\sigma}_e^4}{3} \cdot n^2 \, k_{\max}^3\, (1+ o(1))$.

To show the consistency of our testing statistic under a sequence of local alternatives, we test

equation[equation omitted — 83 chars of source]

where positive sequence $a_n\to 0$ with certain rate while $m(\mathbf{x})\in {\mathcal L}^2$ and ${\mathbb E}[m^2({\bf x})]>0$.

It is pointed out that a sequence of local alternatives covers a wide range of weak endogeneity. The following theorem establishes the consistency of the proposed test under $H_1$.

theo(i) Let the conditions of Theorem (ref) hold. (ii) Let ${\mathbb E}[m^2({\bf x})]>0$. Consider $H_1: \, {\mathbb P}(m_n({\bf x})=a_n\, m({\bf x}))=1$, where $a_n\to 0$ and $a_n^2 \, n \, k_{\max}^{-1/2}\to \infty$ as $n\rightarrow \infty$. We then have under $H_1$, \begin{equation} \frac{\displaystyle1}{\sigma_n} \, T_n\rightarrow_{P} \infty \ \ \ as $n\rightarrow \infty$. \end{equation}

The proofs of Theorems (ref)--(ref) are given in Appendix D. Discussions of Theorems (ref) and (ref) are similar to those for Theorems (ref) and (ref) in Section (ref) above.

Nonlinear and non--separable models

We now show that the proposed SP method can also be extended in a unified way to a number of non-- and semi--parametric regression models associated endogeneity issues, such as

equation[equation omitted — 99 chars of source]

as a semiparametric regression model, along with the following nonparametric regression:

equation[equation omitted — 61 chars of source]

Appendices (ref).1 and (ref).2 below outline the main ideas and steps about how to estimate models ((ref)) and ((ref)) consistently and unbiasedly.

In Appendix (ref).3, we will also discuss one class of binary models of the form:

equation[equation omitted — 94 chars of source]

where the quantities are the same as before.

Semiparametric regression

Observe that model ((ref)) may be motivated as follows:

equation[equation omitted — 162 chars of source]

where a semiparametric projection gives us the following model:

equation[equation omitted — 138 chars of source]

where $\mathbf{x}_v$ is either an observed variable or an IV available to the econometrician satisfying

equation[equation omitted — 80 chars of source]

under which we may identify and estimate $\bm{\alpha}_0$ consistently in the same way as has been done in the relevant literature. Model ((ref)) with ((ref)) has been employed in Example 5.1.

When condition ((ref)) is not satisfied, we consider model ((ref)) for the case where

equation[equation omitted — 189 chars of source]

Note that the proposed approach below is valid regardless of whether ${\mathbb E}[\varepsilon|\mathbf{x}_u]\neq 0$ or ${\mathbb E}[\varepsilon|\mathbf{x}_u]=0$. Let $q_1(\mathbf{x}_v) = {\mathbb E}[y|\mathbf{x}_v]$, $\mathbf{q}_2(\mathbf{x}_v) ={\mathbb E}[\mathbf{x}_u|\mathbf{x}_v]$, $\widetilde{\mathbf{x}}_u=\mathbf{x}_u - \mathbf{q}_2(\mathbf{x}_v)$, $\widetilde{y}=y - q_1(\mathbf{x}_v)$, $q_3(\mathbf{x}_v) = {\mathbb E}[\varepsilon|\mathbf{x}_v]$, $\eta = \varepsilon - q_3(\mathbf{x}_v)$, $m(\widetilde{\mathbf{x}}_u) = {\mathbb E}[\eta|\widetilde{\mathbf{x}}_u]$ and $e=\eta - m(\widetilde{\mathbf{x}}_u)$. Then model ((ref)) can be written as

equation[equation omitted — 250 chars of source]

with ${\mathbb E}[e|\widetilde{\mathbf{x}}_u]=0$. Notationally, model ((ref)) is the same as model (2.3) with $(y, \mathbf{x})$ being replaced by $(\widetilde{y}, \widetilde{\mathbf{x}}_u)$, respectively. The main parameter--of-interest, $\bm{\alpha}_0$, can then be estimated in the same way as in Section 2.1 based on a sampling version of the form:

equation[equation omitted — 126 chars of source]

where $\widetilde{y} = y - \widehat{q}_1(\mathbf{x}_v)$ and $\widetilde{\mathbf{x}}_u = \mathbf{x}_u - \widehat{\mathbf{q}}_2(\mathbf{x}_v)$, in which $\widehat{q}_1(\cdot)$ and $\widehat{\mathbf{q}}_2(\cdot)$ can be constructed as in Section 2.1, before $m(\cdot)$ can be estimated in the same way as in Section 2.1.

It is noted that our discussion covers the case where $\mathbf{x}_u$ is a binary variable, and $\mathbf{x}_v$ reduces to a univariate fixed--design variable of the form $\mathbf{x}_v = \tau \in [0,1]$.

Nonparametric regression

Consider model ((ref)), where $g(\mathbf{x})$ is an unknown function of interest, and the error term $\varepsilon$ satisfies ${\mathbb E}[\varepsilon]=0$, but ${\mathbb E}[\varepsilon|\mathbf{x}] \neq 0$. In such cases, there are several different approaches proposed in dealing with potential endogeneity issues. One of the existing methods is the so--called “Nonparametric IV" approach, as discussed and reviewed in a recent paper by cff2025, for example, which can be summarized as follows.

Let $\mathbf{z}$ be an IV such that ${\mathbb E}[\varepsilon|\mathbf{z}]=0$. One may then use ((ref)) to derive

equation[equation omitted — 195 chars of source]

for which $g_1(\mathbf{z})={\mathbb E}[y|\mathbf{z}]$ can be estimated by an existing nonparametric method before one might be able to recover $g(\mathbf{x})$ by solving an inverse problem ${\mathbb E}[g(\mathbf{x})|\mathbf{z}]={g}_1({\bf z})$.

As stressed before, the main drawback is that such an approach depends on the availability and validity of $\mathbf{z}$ as the IV, in addition to addressing ill--posed inverse issues.

Our discussion is to offer an alternative by decomposing $\varepsilon$ and rewriting model ((ref)) as

equation[equation omitted — 163 chars of source]

where ${\mathbb E}[e|\mathbf{x}]=0$ is automatically satisfied.

We first consider the case where $m(\mathbf{x}) = {\mathbb E}[\varepsilon|\mathbf{x}] = m(\mathbf{x}, \bm{\gamma}_0)$. Model ((ref)) then becomes

equation[equation omitted — 154 chars of source]

which is the same notation as discussed in Appendix A.2.1, where $\widetilde{g}_1(\mathbf{x}) = g(\mathbf{x}) - \mu_0$ with $\mu_0 = {\mathbb E}[g(\mathbf{x})]$. Both the identifiability and estimation of $(\mu_0, \bm{\gamma}_0, \widetilde{g}_1(\cdot))$ follows from that for $(\bm{\theta}_0, m(\cdot))$ discussed in Appendix A.2.1.

It is pointed out that model ((ref)) allows the case where $g(\mathbf{x})$ is nonlinear in $\mathbf{x}$ and $m(\mathbf{x}, \bm{\gamma}_0) = \mathbf{x}^{\top}\bm{\gamma}_0$. Consequently, the case where $(\varepsilon, \mathbf{x})$ follows a joint Gaussian distribution can be covered in ((ref)).

In the case where $m(\mathbf{x})={\mathbb E}[\varepsilon|\mathbf{x}]$ is specified nonparametrically, substantially new developments are required. We therefore wish to write up such details into a different paper.

Binary regression

Regression models with binary dependent variables have many theoretical investigations and empirical studies, and the relevant literature is comprehensive, such as lewbel1998, and blundell2004, for example. To show that the proposed SP method is also useful to address certain types of endogeneity involved in binary models, we start with the following linear probability model:

equation[equation omitted — 104 chars of source]

where the exogeneity case of ${\mathbb E}[\varepsilon|\mathbf{x}]=0$ has been discussed in the relevant literature, such as Chapter 11 of jsmw2018. \, We then rewrite model ((ref)) as

equation[equation omitted — 88 chars of source]

where $m(\mathbf{x}) = {\mathbb E}[\varepsilon|\mathbf{x}]\neq 0$, $e = \varepsilon - m(\mathbf{x})$, and ${\mathbb E}[e|\mathbf{x}]=0$ by definition. Model ((ref)) can then be estimated in the same way as in Section 3, although $y$ is now a binary response variable.

Meanwhile, we consider a non--separable binary model of the form:

equation[equation omitted — 164 chars of source]

which can be rewritten as ${\mathbb E}[y|\mathbf{x}] = {\rm P}_{e}(e\leq \mathbf{x}^{\top} \bm{\theta}_0 - m(\mathbf{x})) = F_{e}(\mathbf{x}^{\top} \bm{\theta}_0 - \mathbf{v}(\mathbf{x})^{\top} \bm{\gamma}) = F_{e}(\mathbf{w}_{-}(\mathbf{x})^{\top} \bm{\theta}_{-})$ when ignoring the approximation error term and $m(\mathbf{x})$ is replaced by $\mathbf{v}(\mathbf{x})^{\top} \, \bm{\gamma}$ as in Section 2, where $F_{e}(\cdot)$ denotes the cumulative distribution function (CDF) of $e$, $(\mathbf{v}(\mathbf{x}), \bm{\gamma})$ is the same as in Section 2, $\mathbf{w}_{-}(\mathbf{x}) = \left(\mathbf{x}^{\top}, -\mathbf{v}^{\top}(\mathbf{x})\right)^{\top}$ and $\bm{\theta}_{-} = \left(\bm{\theta}_0^{\top}, \bm{\gamma}^{\top}\right)^{\top}$.

When $F_{e}(\cdot)$ is parametrically known, the parameter--of--interest, $\bm{\theta}_0$, can then be estimated consistently and unbiasedly by MLE. When $F_e(\cdot)$ is nonparametrically unknown, it can be expanded by an infinite sum of known CDFs before a semiparametric MLE estimation method may be developed.

Discussion of other estimation methods

SP estimation efficiency

Consider the case of $d=1$ and $\sigma_i^2 \equiv \sigma_e^2$ for notational simplicity in the following derivations. Recall the standard OLS estimator by

equation[equation omitted — 226 chars of source]

where $s_n \equiv: \sum_{i=1}^n \mathbf{x}_i \, \mathbf{x}_i^{\top}$.

Let $\bm{\Sigma}_{xm} = {\mathbb E}[e^2] \, {\mathbb E}\left(\left[\mathbf{x} \, \mathbf{x}^{\top}\right]\right) + {\mathbb E}\left(\left[\mathbf{x} \, m(\mathbf{x}) - {\ell}_{\rm endo}\right] \, \left[\mathbf{x} \, m(\mathbf{x}) - {\ell}_{\rm endo}\right]^{\top}\right)$, $\bm{\Sigma}_{\rm LS} = {\mathbb E}^{-2} \left[\mathbf{x}_1^2\right]\, \bm{\Sigma}_{xm}$ and ${\ell}_{\rm endo} = {\mathbb E}[\mathbf{x} \, m(\mathbf{x})]$, and $\bm{\Sigma}_{\rm SP} = \sigma_e^2 \, \bm{\Sigma}_x^{-1}$ with $\bm{\Sigma}_x = {\mathbb E}[\mathbf{x}_1^2] - \sum_{j=1}^{\infty} {\mathbb E}^2\left[\mathbf{x}_1 \, \psi_j(\mathbf{x}_1)\right]$. We now show that $\bm{\Sigma}_{\rm SP} \leq \bm{\Sigma}_{\rm LS}$ as follows:

eqnarray[eqnarray omitted — 628 chars of source]

when ${\mathbb E}\left(\left[(\mathbf{x}_1 \, m(\mathbf{x}_1) - {\ell}_{\rm endo})\right]^2\right) \geq \frac{\bm{\sigma}_{11} \, \bm{\sigma}_{12}}{\bm{\sigma}_{11} - \bm{\sigma}_{12}} \, \sigma_e^2$, where $\bm{\sigma}_{11} = {\mathbb E}[\mathbf{x}_1^2]$ and $\bm{\sigma}_{12} = \sum_{j=1}^{\infty} {\mathbb E}^2\left[\mathbf{x}_1 \, \psi_j(\mathbf{x}_1)\right]$.

Equation ((ref)) remains true even in the case of ${\ell}_{\rm endo}=0$ and $\bm{\sigma}_{12} = 0$. In other words, $\widehat{\bm{\beta}}_{\rm SP} \equiv: \widehat{\bm{\beta}}$ is more efficient than $\widehat{\bm{\beta}}_{\rm LS}$ even when ${\ell}_{\rm endo}=0$.

GMM estimation method

Recall the notation and symbols introduced in Section 3. Letting $\mathbf{Q}_{d\times d}$ be a known positive definite weight matrix to be chosen by the user, we estimate $\bm{\beta}_0$ by

align*[align* omitted — 314 chars of source]

which offers a closed--form expression as follows:

align*[align* omitted — 360 chars of source]

It is known from the discussion in Section 3 that

align*[align* omitted — 156 chars of source]

We then have

eqnarray[eqnarray omitted — 1,181 chars of source]

The second term is $o_P(1)$ as $k\to \infty$ under Assumption 3.1(iii)(iv) as shown in the proof of Theorem 3.1. For the first term, the conditional covariance matrix is

eqnarray[eqnarray omitted — 734 chars of source]

as $(n,k)\to(\infty, \infty)$ under the conditions of Theorem 3.1(ii) since ${\bf Q}$ and $\bm{\Sigma}_x(k)$ are invertible. The result is the same as Theorem 3.1.

Estimation of weakly identified linear models

We start with the linear case of $m(\mathbf{x}) = \left(\mathbf{x} - {\mathbb E}[\mathbf{x}]\right)^{\top} \bm{\gamma}_0$. Let us define a truncated version of $\bm{\Sigma}_x$ of the form:

equation[equation omitted — 243 chars of source]

If we assume that we can expand $\mathbf{x} - {\mathbb E}[\mathbf{x}]= \sum_{j=1}^{\infty} \psi_j(\mathbf{x}) \, \bm{\gamma}_j$ by the same orthonormal series: $\{\psi_j(\cdot): j\geq 1\}$ as used in Section 2, we have each given $k\geq 1$

equation[equation omitted — 180 chars of source]

which has a reduced rank when ${\mathbb E}[\mathbf{x}]\neq 0$ and $d>1$. In this case, the proposed LASSO selection estimation method in Section 5.1, which allows for $m(\mathbf{x})$ to be linear in $\mathbf{x}$, addresses such reduced--rank issues. The finite--sample evaluation results in Section 5 and Appendix (ref) support the LASSO selection method.

In the case of $d=1$ and ${\mathbb E}[x]\neq 0$, it can be seen that the SP method itself is directly applicable, and Theorem 3.1 remains true. Example B.2.2 of Appendix (ref) shows that without using the LASSO selection method, the SP selection method works well with commonly used functions, including $m({x}) = \left({x} - {\mathbb E}[{x}]\right) {\gamma}_0$ and second--order polynomial functions.

When ${\mathbb E}[\mathbf{x}] = 0$, model (2.3) reduces to

equation[equation omitted — 223 chars of source]

which means that one may only be able to correctly identify $\left(\bm{\beta}_0 + \bm{\gamma}_0\right) = {\mathbb E}^{-1}\left[\mathbf{x} \, \mathbf{x}^{\top}\right] \, {\mathbb E}\left[\mathbf{x} \, y\right]$ collectively, rather than $\bm{\beta}_0$ individually.

For model ((ref)) in the case of ${\mathbb E}[\mathbf{x}] = 0$, meanwhile, equation ((ref)) implies that as $k\rightarrow \infty$

equation[equation omitted — 129 chars of source]

Due to ((ref)), for model ((ref)), we cannot assume Assumption 2.1(iii). Instead we replace it by {\bf Assumption 2.1(iii)*}: Assume that there are a positive definite matrix of real numbers, $\mathbf{Q}_n$, and another positive definite matrix, $\bm{\Sigma}_{x, \ast}$, such that $\lambda_{\rm \min}(\mathbf{Q}_n) \rightarrow \infty$, $ n \, \lambda_{\rm \min}(\mathbf{Q}_n^{-1}) \rightarrow \infty$ and $\mathbf{Q}_n \, \bm{\Sigma}_x(k)\rightarrow_P \bm{\Sigma}_{x, \ast}$ as $n\rightarrow \infty$, where $\lambda_{\min}(A)$ denotes the smallest eigenvalue of matrix $A$.

In view of the proof of Theorem 3.1, replacing Assumption 2.1(iii) by Assumption 2.1(iii)*, it can then be established that the following asymptotic normality holds:

equation[equation omitted — 371 chars of source]

which still offers a consistent estimator for $\bm{\beta}_0$, with a reduced rate of convergence of an order of $\left(\sqrt{n} \, \mathbf{Q}_n^{-1/2}\right)$, slower than $\sqrt{n}$, in which

equation[equation omitted — 295 chars of source]

Finite--Sample Evaluations

Implementational Issues

We discuss several important issues about how to implement the proposed SP estimation procedure in practice. We offer our recommendations on the choice of orthonormal series functions and the truncation parameters involved in the proposed SP estimation.

Dimension reduction

The main model and estimation method proposed in Section 3 remains valid for the case where the dimensionality, $d$, is large but fixed in theory, although we assume in Section 2 that we focus on the case where $d$ is small. To explain this in a bit more detail, we recall from Section 2 that we expand $m({\bf x})$ as $m(\mathbf{x}) = \sum_{j=1}^{\infty} \psi_j(\mathbf{x}) \, \gamma_j$, in which the dimensionality of $\mathbf{x}$ is only involved in the chosen series $\{\psi_j(\mathbf{x}): j\geq 1\}$. For ease of implementation, we suggest using $\psi_j(\mathbf{x}) = \prod_{k=1}^d \psi_{jk}(x_k)$ when $d\geq 2$, where $\{\psi_{jk}(\cdot): k\geq 1\}$ is an array of univariate series functions, and $\mathbf{x} = (x_1, \cdots, x_d)$.

Meanwhile, some other dimension reductions might be employed, such as an additivity structure of the form: $m(\mathbf{x}) = \sum_{j=1}^d m_j(x_j)$ directly as in ln1995, in which it is expected that the construction of $\mathcal{S}$ should be a direct sum of $\mathcal{S}_j$ for $1\leq j\leq d$. There are also dimension reduction methods proposed by involving certain types of single--index or multi--index modelling methods (see, dg2025, for proposing an additive single--index structure form for each $\psi_j(\cdot)$). Our experience with the finite--sample studies in Section 5 and Appendix (ref) shows that the multiplicative form: $\psi_j(\mathbf{x}) = \prod_{k=1}^d \psi_{jk}(x_k)$ works well and better than some other competing forms available in the relevant literature.

Choice of series functions

In both theory and practice, we need not require orthogonality on $\{\psi_l(\cdot): l\geq 1\}$, although the orthogonality assumption simplifies the notation involved in the theoretical derivations. In practice, our experience suggests using either one of the following series, or a mixture of both.

1. \, The probabilist's Hermite polynomials $\{H_j(x),j\ge 1\}$ are a set of orthonormal basis functions defined on $L^2(\mathbb{R}, \exp(-x^2/2))$ for the univariate setting;

2. \, Let the trigonometric polynomials be $p_j(x) = \sqrt{2}\cos(\pi j x)$ with $j\ge 1$. Lemma E.1 of Appendix E in the online supplement shows that $\{p_j(x): j\ge 1\}$ is an orthogonal set of basis functions on $L^2([a, b])$ as long as $(a,b)$ are different integers. For the case $a=0$ and $b=1$, $\{p_j(x): j \ge 1\}$ constitutes an orthonormal family.

3. \, In the multivariate setting, we propose using the following form of either $p_j(\mathbf{x}) = \prod_{k=1}^d p_{jk}(x_k)$ or $H_j(\mathbf{x}) = \prod_{k=1}^d h_{jk}(x_k)$ in simulations and empirical applications, where $p_{jk}(\cdot)$ and $h_{jk}(\cdot)$ are the corresponding univariate functions that may be chosen as in Steps 1 and 2 above.

As discussed in Section 4.2.3 and Appendix C.4 of dg2025, we need not know the distributional structure of the data under analysis in practice as long as the support of $\{\mathbf{x}_i: i\geq 1\}$ becomes available to the practitioner.

In theory, the choice of $\{\psi_j(\cdot): j\geq 1\}$ can be flexible as long as $\bm{\Sigma}_x>{\bf 0}$. As discussed in Section 5.1 of the main submission, moreover, the proposed LASSO selection not only facilitates the choice of an optimal orthonormal series for $\{\psi_j(\cdot): j\geq 1\}$ in practice, but also helps address possible reduced--rank issues as alluded in Appendix (ref) above.

Main simulations

Finite--sample properties of the estimation theory

In this section, we consider a number of scenarios to demonstrate the finite-sample performance of the proposed SP estimation method under orthonormality. To show that the orthonormality and even the orthogonality on $\{\psi_j(\cdot): j\geq 1\}$ may all be relaxed, we present extensive numerical evaluations in Appendix E of the online supplementary document to demonstrate that the SP method still works well numerically. We provide the code at \url{https://github.com/pengbin430/SIV_2026/tree/main}.

{\bf Example B.2.1}: We consider the following data generating processes:

itemize• Case A: $y_i = \alpha_0 +x_i \, \beta_0 + \varepsilon_i$, $(\alpha_0,\beta_0) =(1,1)$, $\varepsilon_i = m(x_i) + e_i$, $e_i\sim N(0,1)$, $x_i \sim U(0, \pi)$ and $m(x)= 2.5\cos(3 \, x)+0.5\cos(5 \, x)$; • Case B: Consider Case A, but choose $m(x)$ as $m(x)=\sum_{j=1}^{\infty} \gamma_j\cos(j x)$, and $\{ y_i, x_i\}$ are observable, in which $\gamma_j =4\cdot (0.9)^j$ if $j\in \{4(\ell-1)+3\mid \ell=1,2,\ldots \}$, and $\gamma_j =0$, otherwise; • Case C: $y_i =x_i^2 \, \beta_0 +\varepsilon_i$, $x_i\sim U(0, \pi)$ and $m(x_i)=12\pi (x_i-\frac{\pi }{2})$; • Case D: Consider Case A with one additional exogenous regressor $y_i = \alpha_0 +x_i \, \beta_0+w_i\beta_1 + \varepsilon_i$, where $w_i\sim N(1,1)$ and $\beta_1=1$.

In Appendix (ref).3 below, we show that $E[x_im(x_i)]\ne 0$ and $E[m(x_i)]=0$ for all cases, so the endogeneity exits. Case C can also be considered as a special case of the NLS framework discussed in Appendix A.2. Both Cases B and C have truncation residuals. As shown in the justification, Case B has more sparsity in $\{\gamma_j \}$, so it is reasonable to expect our approach has better finite sample performance for Case B.

After identifying $\phi_j(x)$'s with nonzero coefficients from the LASSO method, we then construct $\{\psi_j(\cdot): j\geq 1\}$ to run SP to obtain $\widehat{\beta}_{\rm SP}$. We alway enforce $\mathbf{V}_k(\cdot)$ to include the constant term 1 in order to remove the intercept.

{

table[table omitted — 2,324 chars of source]

}

We repeat the above procedure $R=1000$ times and then report the following measures in Table (ref):

eqnarray[eqnarray omitted — 186 chars of source]

where $\overline{\beta} = \frac{1}{R}\sum_{r=1}^R\widehat{\beta}_r$ with $\widehat{\beta}_r \in\{\widehat{\beta}_{\rm LS},\widehat{\beta}_{\rm SP}\}$ in each replication. Additionally, we calculate $\left|\frac{\Delta \beta_{SP}}{\Delta \beta_{LS}}\right|$ in Table (ref) in order to show the relative magnitude of the biases associated with both OLS and SP methods. As indicated in Table (ref), the OLS method is apparently biased and has large standard deviation. The SP method works well for all cases. Although the magnitudes of biases vary in different cases, the ratio $\left|\frac{\Delta \beta_{SP}}{\Delta \beta_{LS}}\right|$ is reasonably stable across all cases. The only exception is the ratio associated with $\beta_1$ of Case D, which should be expected due to the fact that $w_i$ is an exogenous variable.

{

table[table omitted — 3,022 chars of source]

}

We then examine the LASSO selection method to further reiterate our discussion in the above sections. Specifically, in Table (ref), we report the following measure: $P_j =\frac{1}{R}\sum_{r=1}^R I(\phi_j\in S_r)$, where $S_r$ stands for the set $S$ selected by the LASSO method in the $r^{th}$ replication. Therefore, $P_j$ measures the probability of $\phi_j$ being selected over $R$ replications.

It is clear that that the LASSO method works well. A few findings emerge. Take Case C as an example. Justification of Example B.2.1 given in Appendix (ref) shows that the odd indexed $\phi_j$'s have nonzero coefficients, while the even indexed $\phi_j$'s have coefficients 0. Consequently, Table (ref) shows that the odd indexed $P_j$'s are often selected with high probability, while the even indexed $P_j$'s are barely selected.

We then add some simulation results for the case where $m(\mathbf{x})$ itself contains a linear component in each case under examination in Example B.2.2 below.

{\bf Example B.2.2}: We now consider model $y_i = x_i \, \beta_0 + \varepsilon_i$, $\varepsilon_i = m(x_i) + e_i$, $e_i\sim N(0,1)$ and $x_i \sim U(0, \pi)$ under the following scenarios:

itemize• Case E: $m(x) = 1.5 \, (\cos(x) + 2 \,x - \pi) $; • Case F: $m(x) = 1 + (2 \, x - \pi) - \frac{3}{\pi^2} \, x^2$; and • Case G: $m(x) = 1.5 \, (x - 0.5 \, \pi)$.

Apparently, the restrictions: ${\mathbb E}[m(x_i)]=0$ and ${\mathbb E}[x_i \, m(x_i)]\neq 0$ can be verified in a similar way to those for Example B.2.1. We summarize the simulation results in Table (ref) below, wherein the biases and standard deviations (in parentheses) are calculated in exactly the same way as before.

table[table omitted — 625 chars of source]

Table (ref) obviously shows that the SP method works well numerically. Moreover, the numerical results are better than those in Example B.2.1, probably because the elementary functional forms of $m(\cdot)$ considered in Cases E--G were more accurately approximated by the trigonometric series $\left\{\psi_j(x) = \sqrt{\frac{2}{\pi}} \, \cos(j\, x): j\geq 1\right\}$ for the case of $x\in U[0, \pi]$.

Finite--sample properties of the testing theory

In this section, we evaluate the finite--sample property of the test proposed in Section 3.2.

{\bf Example B.2.3}: Without loss of generality, we consider Cases A and B of Example B.2.1, and slightly make the following modification in order to evaluate size and power respectively.

eqnarray*[eqnarray* omitted — 123 chars of source]

where $x_i \sim U(0, \pi)$, and $a_{n,j} =c_j \, \sqrt{\frac{k_{\max}}{n}}$ for $1\leq j\leq 4$. Without of loss generality, we chose $c_1\equiv 0$ to evaluate the sizes, and $c_2 \equiv 0.75$, $c_3 \equiv 1$, and $c_4 \equiv 1.25$ to evaluate the local powers. Accordingly, we have $\{\psi_{j}(x)=\sqrt{2/\pi}\cos(jx)\mid j \ge 1 \}$, which are defined precisely in multiple places. When calculating the test statistic, we chose $k_{\min}=\lceil \log \log n\rceil$ and $k_{\max}=\lceil 3 \log \log n\rceil$. By construction, we have

eqnarray[eqnarray omitted — 89 chars of source]

so the requirement to ensure local power properties is fulfilled. We generate $m(x)$ in the following two forms, which are almost identical to those in Example B.2.1.

itemize• Case A: $m(x)= 4\cos(3 x)+4\cos(5 x)$; • Case B: $m(x)=\sum_{j=1}^{\infty} \gamma_j\cos(j x)$ in which $\gamma_j =4\cdot (0.9)^j$ if $j\in \{4(\ell-1)+3\mid \ell=1,2,\ldots\}$; $\gamma_j =0$, otherwise.

To improve the finite sample performance of the size function in each case, we propose using the following wild bootstrap procedure (see, for example, gg2008) to obtain bootstrapping critical values.

enumerate• Obtain $\widetilde{e}_i$ as specified under (3.5), and generate bootstrap sample $\{y_i^* \}$, where $y_i^* \equiv x_i \, \widehat{\beta}+ \widetilde{e}_i \eta_i$ and $\{\eta_i\}$ are i.i.d. sample generated from $N(0,1)$. • Given $\{y_i^*, x_i\}$, we run regression to obtain $\{\widetilde{e}_i^*\}$ as in Step 1, and calculate the corresponding bootstrap statistic $L_n^*$. • We repeat the above steps, say, 400 times to obtain the 95% coverage set $\mathscr{C}_n$.

After $R=1000$ simulation replications, we report the following rejection rate for $n=200,300, 400$ respectively: $\text{RJ} =\frac{1}{R} \sum_{r=1}^R I(L_{n,r}\not \in \mathscr{C}_{n,r}),$ where $L_{n,r}$ and $\mathscr{C}_{n,r}$ respectively stand for the values of $L_n$ and $\mathscr{C}_n$ obtained in the $r^{th}$ replication. The results are summarized in Table (ref).

table[table omitted — 563 chars of source]

In Table (ref), the rejection rates with $c_1=0$ are around the nominal rejection rate (i.e., 5%). When using $c_2$, $c_3$ and $c_4$, the rejection rates reflect the local power of the proposed test. It is not surprising that the local power varies across two cases. In Case B, the rejection rates are slightly lower than those in Case A, which might be due to the impact of the truncation residual. In summary, Table (ref) shows that the proposed test can detect the weakest possible endogeneity at an optimal order of $n^{-1/2} \, \delta_n$ for such $\delta_n$ that diverges to $\infty$ at one of the slowest possible rates of an order of $\log(\log(n))$.

Verification of the simulation designs in Examples B.2.1 and B.2.2

We consider Case C of Example B.2.1 only here. As the derivations for Cases A, B and D of Example B.2.1 are similar and simpler, and for Example B.2.2 are also similar, we omit the details. To show that $E[m(x_i)]=0$ and $E[x_im(x_i)]\ne 0$, it is sufficient to consider

eqnarray*[eqnarray* omitted — 281 chars of source]

so endogeneity exists.

We then show that

eqnarray*[eqnarray* omitted — 546 chars of source]

and

eqnarray*[eqnarray* omitted — 413 chars of source]

Therefore, we have shown that $m(x)$ can be fully expanded by the odd terms of $\{\cos(j \, x)\}$, while fully expanding the regressor $x^2$ requires all of $\{\cos(j \, x)\}$.

Extra empirical analysis

{\bf Example 5.1} (continued)

As alluded before, we change the basis functions to $\phi_j(x) =(x-\pi/2)^j$ for $j\geq 1$, and keep everything else unchanged. Obviously, the new set of basis functions are not even orthogonal to each other in the space $L^2([0,\pi])$. The LASSO method selects $\phi_2(x) =(x-\pi/2)^2$ only, so $\widehat{m}(x)= -0.0643(x-\pi/2)^2$.

We then obtain the following table, wherein the last two rows are corresponding to the results obtained from the 2sLS method by using $(z_1,\phi_2)$ and $\phi_2(x)=(x-\pi/2)^2$ as IVs, respectively.

table[table omitted — 454 chars of source]

It is observed that the new estimation results are very similar to those reported in Section 5. This may indicate the robustness of our method, and the choice of $\{\phi_j(\cdot): j\geq 1\}$ and whether it is orthonormal or not may not affect the corresponding estimation results very much.

Possible sources of endogeneity

As an additional evidence to show that there is some substantial nonlinear endogeneity involved in Example 5.1, we plot the estimated $m(x)$ with its 95% confidence interval (obtained via wild bootstrap) in Figure (ref). It is clear the SP method supports a clearly nonlinear functional form for the specification of $m(x)$.

{\samepage {

figure[figure omitted — 128 chars of source]

}}

As briefly pointed out in the previous sections of this paper, an additional novelty of the SP approach is that we are able to estimate $\varepsilon$ of model (2.1) by $\widehat{\varepsilon} = y - \mathbf{x}^{\top} \widehat{\boldsymbol{\beta}}$. As a consequence, we may be able to make use of $\widehat{\varepsilon}$ to establish an approximate regression model of the form:

equation[equation omitted — 110 chars of source]

for the econometrician to try to reveal possible sources of endogeneity, where each $\mathbf{w}_i$ is a vector of variables possibly containing omitted variables and/or errors in variables, which may be highly correlated with $\mathbf{x}$, and ${\eta}$ is the error term involved.

Due to the fact that $\varepsilon$ is usually unobservable and latent, it is reasonable to assume the non--separable functional form in ((ref)). If each $\mathbf{w}_i$ is available to the econometrician, the discussion in Appendices A.2 and (ref) suggests that one may be able to reveal possible sources of endogeneity by estimating the functional form of $G(\mathbf{w}; \cdot)$ with respect to $\mathbf{w}$.

Let us now come back to Example 5.1. We would like to investigate how $\varepsilon$ might be uncorrelated with $\{z_j: 1\leq j\leq 5\}$ used above. We compute the sample averages of the following quantities: $E[\widetilde{x}\cdot m(\widetilde{x})]$ and $E[z_j\cdot \varepsilon]$, where $\widetilde{x}$, $m(\widetilde{x})$ are defined in (5.3), $z_j$'s are the different IV variables adopted above, and those unobservable (i.e., $m(\cdot)$ and $\varepsilon = y- \widetilde{\alpha}_0- \widetilde{x}\, \widetilde{\beta}_0$ are replaced with their estimates via the SP method).

As shown in Table (ref), $E[\widetilde{x}\cdot m(\widetilde{x})]$ is clearly non-zero, which justifies the existence of the endogeneity. As an valid IV, we would expect $E[z_j\cdot \varepsilon]$ all to be zero. However, only $z_1$ and $z_5$ generate averages sufficiently closed to 0. Numerically, it offers a reason to explain why the results associated with $z_1$ and $(z_1, z_5)$ in Table 2 are close to that from the SP approach. However, the other quantities: $E[z_j\cdot \varepsilon]$, for $2\leq j\leq 4$, don't justify the validity of those IVs.

table[table omitted — 495 chars of source]

Finally, we estimate the following three nonparametric regressions models:

equation[equation omitted — 198 chars of source]

where $\widetilde{x}$, $z_1$ and $\varepsilon$ have been defined previously, $g_1(\cdot)$, $g_2(\cdot)$ and $g_3(\cdot)$ measure possible relationships among them, and $\xi_1$, $\xi_2$ and $\xi_3$ are the error terms.

{\samepage

figure[figure omitted — 147 chars of source]

}

The nonparametrically estimated relationships, along with their corresponding $95\%$ confidence intervals, are plotted in Figure (ref). The second plot clearly shows that the relationship between $\varepsilon$ and $z_1$ is clearly nonlinear, although being symmetrically fluctuating around zero. For simplicity, we use the GCV method of gao2002 to select the truncation parameter here. In the second case, the GCV suggests the number of truncation parameters being $20$. The second plot infers that the typical IV method might be modelled beyond linear regression.

The first and third plots together infer the relationship between $x$ and $z_1$ has a linear trend but not exactly linear, as the truncation parameters for these two regressions are 13 and 7 respectively according to the GCV approach. This further explains why the results associated the SP method and the IV method (using $z_1$) are close to each other in Table 2.

In summary, due to the unobservable nature of $\varepsilon$, it is difficult to establish and reveal a relationship between $\varepsilon$ and potential instrumental variables. Without assuming the knowledge of such relationships, the existence and the validity of potential IVs, the proposed SP approach offers a feasible and robust way to deal with certain types of endogeneity issues.

As discussed above, the proposed SP estimation method is applicable to time--series data. So we now have a look at the following application.

Endogeneity in stock return

We apply the SP method to revisit the classic market model (M1997, campbell1998econometrics). See Eq. (4.3.2) of campbell1998econometrics for example.

Consider the following linear model:

eqnarray[eqnarray omitted — 67 chars of source]

where $\varepsilon_t=m(x_t) + e_t$, $y_t$ represents daily stock return for a chosen stock, and $x_t$ stands for the market portfolio. In this study, we use daily return of the SP500 index as $x_t$. We aim to argue that endogeneity should also be taken into account for the market model.

For the response $y_t$, we collect data for Apple, Google, Meta and Microsoft from Yahoo finance covering period from the first trading day of 2012 to the last trading day of 2020, which gives $n=2169$ observations for us to conduct regressions. As in our study about the return to schooling, we normalize $x_t$ to get $\widetilde{x}_t$, so the latter belongs to $[0,\pi]$. Thus, the model used for regression is again written as

eqnarray[eqnarray omitted — 98 chars of source]

For each stock, we consider (i) The SP method; \ (ii) The LS method; and (iii) The 2sLS method using the basis functions selected by the SP method as IVs.

{

figure[figure omitted — 133 chars of source]

}

The estimated $m(x)$ for different companies are given as follows:

enumerate• Apple: $\widehat{m}(x) = 0.00348\phi_2(x)$; • Google: $\widehat{m}(x) = 0.03176\phi_1(x) +0.00750\phi_3(x)-0.00187\phi_4(x)$; • Meta: $\widehat{m}(x)= 0.02730\phi_1(x) -0.00647\phi_5(x) $; • Microsoft: $\widehat{m}(x) = 0.00352\phi_2(x) -0.00123\phi_4(x)$.

For the purpose of demonstration, we also plot $\widehat{m}(x)$ from Google with its 95% confidence interval in Figure (ref) below, which clearly shows some nonlinear fluctuation of the estimated endogeneity component represented by $\widehat{m}(\cdot)$.

table[table omitted — 1,350 chars of source]
table[table omitted — 378 chars of source]

Tables (ref) and (ref) summarize the results, revealing several key findings. First, Table (ref) indicates similar features about the SP estimates to those discussed in Section 5 with relatively smaller RMSEs in comparison with those of LS and 2sLS estimates. Second, Table (ref) confirms the presence of endogeneity, as the estimated values of $E[\widetilde{x} \, m(\widetilde{x})]$ are negative across all companies, which should be expected given that these four companies account for approximately more than 19% of the entire S&P 500 index.

Third, the estimated $\beta_0$ for all companies is greater than 1 indicating that these four stocks are more volatile than the market after accounting for the endogeneity.

Bias correction for weak endogeneity

This section discusses about how to eliminate the bias term by a simple jackknife method.

To show the main idea, we focus on a weak endogeneity case of the form:

equation[equation omitted — 159 chars of source]

where $\mu=4$, $x_i\sim U(0,\pi)$, and $e_i\sim N(0,1)$.

We partition our sample in two groups: $\mathscr{N}_1$ and $\mathscr{N}_2$, where $\sharp \mathscr{N}_j =n_j$ for $j=1,2$, and $\mathscr{N}_1\cap \mathscr{N}_2=\emptyset$ before we define the OLS estimators: $\widehat{\bm{\beta}}_{\rm LS, 1} = \left(\sum_{\mathscr{N}_1}\mathbf{x}_i\mathbf{x}_i^\top\right)^{-1} \sum_{\mathscr{N}_1}\mathbf{x}_i y_i$ and $\widehat{\bm{\beta}}_{\rm LS, 2} = \left(\sum_{\mathscr{N}_2}\mathbf{x}_i\mathbf{x}_i^\top\right)^{-1} \sum_{\mathscr{N}_2}\mathbf{x}_i y_i$.

Define $\widehat{\bm{\beta}}_{\rm BC} = c \, \widehat{\bm{\beta}}_{\rm LS, 1} + (1-c) \, \widehat{\bm{\beta}}_{\rm LS, 2}$. Let $n=n_1 + n_2$. Simple algebra shows that

eqnarray[eqnarray omitted — 565 chars of source]

Note that ${\rm Var}\left(\frac{c}{\sqrt{n_1}}\sum_{\mathscr{N}_1}\mathbf{x}_i e_i+ \frac{(1-c)\sqrt{n_1}}{n_2}\sum_{\mathscr{N}_2}\mathbf{x}_i e_i\right) = 2 \, c^2\, {\mathbb E}\left[\mathbf{x} \, \mathbf{x}^{\top}\right] \sigma_e^2$ due to requiring $c = (c-1)\frac{\sqrt{n_1}}{\sqrt{n_2}}$. Note also that while the choice of $c$ affects the variance component, it doesn't have any impact on the bias evaluation as long as $c>1$ satisfies $c = (c-1)\frac{\sqrt{n_1}}{\sqrt{n_2}}$.

We therefore choose $c=1.5$ and then $n_2=\left[\frac{n_1}{9}\right]$ in the following numerical evaluation. We consider the case of $n_1\in \{270, 405, 540\}$, so the value of $n_2$ is defined accordingly as follows. We report the absolute bias and the standard derivation based $R=1000$ replications as follows:

eqnarray[eqnarray omitted — 204 chars of source]

where $\overline{\beta} =\frac{1}{R}\sum_{r=1}^R \widehat{\beta}_r $, and $\widehat{\beta}_r$ stands for the value of $\widehat{\beta}_{\rm LS}$ or $\widehat{\beta}_{\rm BC}$ in the $r^{th}$ replication.

table[table omitted — 479 chars of source]

The results are summarized in Table (ref). It is clear that $\widehat{\beta}_{\rm LS}$ yields a much larger bias in each individual case, although $\widehat{\beta}_{\rm LS}$ has a smaller standard derivation correspondingly. The unbiased estimate $\widehat{\beta}_{\rm BC}$ has been obtained by partitioning our full sample size $n=n_1+n_2$, as a consequence, it sacrifices the standard deviations slightly.

Proofs for Sections 2, 3 and 5

proof[Proof of Lemma 2.1] (1) Notice that as an element in ${\mathcal S}$, ${\mathcal P}_{\mathcal S}(\xi)=\sum_{j=1}^\infty c_j \psi_j({\bf x})$, while for ${\mathcal P}_{\mathcal S}(\xi)$ to be a projection of $\xi$, $\xi-{\mathcal P}_{\mathcal S}(\xi)$ should have minimum norm among all $\xi-\eta$ with $\eta\in {\mathcal S}$, or equivalently, among all $c_j$. Noting that $\|\xi-{\mathcal P}_{\mathcal S}(\xi)\|^2={\mathbb E}[\xi-{\mathcal P}_{\mathcal S}(\xi)]^2 = {\mathbb E}[\xi^2]-2\sum_{j=1}^\infty c_j {\mathbb E}[\xi\psi_j({\bf x})]+\sum_{j=1}^\infty c_j^2,$ the first order condition gives $c_j={\mathbb E}[\xi\psi_j({\bf x})]$. (2) By Assumption 2.1(iii), for any $\lambda\ne 0$, $\lambda^\top {\bf x}\not\in {\mathcal S}$. Thus, ${\mathcal M}_{\mathcal S}(\lambda^\top {\bf x})=\lambda^\top {\bf x}-{\mathcal P}_{\mathcal S}(\lambda^\top {\bf x})\ne 0$. We then have $${0}< \|{\mathcal M}_{\mathcal S}(\lambda^\top {\bf x})\|^2={\mathbb E}[(\lambda^\top {\bf x}-{\mathcal P}_{\mathcal S}(\lambda^\top {\bf x}))^2] = \lambda^\top \left\{{\mathbb E}[{\bf xx}^\top]-\sum_{j=1}^\infty {\mathbb E}[{\bf x} \psi_j({\bf x})] {\mathbb E}[{\bf x}^\top\psi_j({\bf x})] \right\}\lambda.$$ The assertion holds.
proof[Proof of Lemma 2.2] We now show that $\bm{\beta}_0$ is invariant to the choice of the basis $\{\psi_j(\cdot): j\geq 1\}$ in ${\mathcal S}$. Towards this end, suppose that ${\mathcal S}$ can also be spanned by another orthonormal sequence $\{\widetilde{\psi}_j(\cdot): j\geq 1\}$. Denote two infinite-dimensional vectors $\Psi({\bf x})=(\psi_1({\bf x}), \psi_2({\bf x}), \cdots)^\top$ and $\widetilde{\Psi}({\bf x})=(\widetilde{\psi}_1({\bf x}), \widetilde{\psi}_2({\bf x}), \cdots)^\top$. Thus, $\widetilde{\Psi}({\bf x})=\mathbf{A}\Psi({\bf x}),$ where $A$ is an infinite-dimensional matrix with element $a_{ij}={\mathbb E}[\widetilde{\psi}_j(\mathbf{x})\psi_{i}(\mathbf{x})]$ that is the $i$-th coefficient in the orthogonal expansion of $\widetilde{\psi}_j(\mathbf{x})$ in terms of the basis $\{\psi_j(\cdot): j\geq 1\}$. On the other hand, $a_{ji}$ is also the $j$-th coefficient in the orthogonal expansion of $\psi_i(\mathbf{x})$ in terms of the basis $\{\widetilde{\psi}_j(\mathbf{x}): j\geq 1\}$. Hence, $\Psi({\bf x})=\mathbf{A}^\top \widetilde{\Psi}({\bf x}).$ Since both $\{\widetilde{\psi}_j(\cdot): j\geq 1\}$ and $\{\psi_j(\cdot): j\geq 1\}$ are orthonormal sequences, ${\mathbb E}[\widetilde{\Psi}({\bf x})\widetilde{\Psi}({\bf x})^\top]=\mathbf{I}$ and ${\mathbb E}[\Psi({\bf x})\Psi({\bf x})^\top]=\mathbf{I}$; these imply $\mathbf{A}\mathbf{A}^\top=\mathbf{I}$ and $\mathbf{A}^\top \mathbf{A}=\mathbf{I}$. Therefore, $\mathbf{A}$ is an orthogonal matrix. It follows that \begin{eqnarray} \sum_{j=1}^{\infty} {\mathbb E}[\mathbf{x} \, \widetilde{\psi}_j(\mathbf{x})] \, {\mathbb E}[\mathbf{x}^{\top} \widetilde{\psi}_j(\mathbf{x})] & = &{\mathbb E}[\mathbf{x} \, \widetilde{\Psi}(\mathbf{x})^\top] \, {\mathbb E}[ \widetilde{\Psi}(\mathbf{x})\mathbf{x}^{\top}]={\mathbb E}[\mathbf{x} \, \Psi(\mathbf{x})^\top]\mathbf{A}^\top \mathbf{A} {\mathbb E}[\Psi(\mathbf{x})\mathbf{x}^{\top}] \nonumber\\ & =&{\mathbb E}[\mathbf{x} \, \Psi(\mathbf{x})^\top]{\mathbb E}[\Psi(\mathbf{x})\mathbf{x}^{\top}]=\sum_{i=1}^{\infty} {\mathbb E}[\mathbf{x} \, {\psi}_i(\mathbf{x})] \, {\mathbb E}[\mathbf{x}^{\top} {\psi}_i(\mathbf{x})], \nonumber \end{eqnarray} which implies that the matrix $\bm{\Sigma}_x= {\mathbb E}[\mathbf{x}\, \mathbf{x}^{\top}] - \sum_{j=1}^{\infty} {\mathbb E}[\mathbf{x} \, {\psi}_j(\mathbf{x})] \, {\mathbb E}[\mathbf{x}^{\top} {\psi}_j(\mathbf{x})]$ is invariant to the choice of the basis, so is $\bm{\Sigma}_{xy} = {\mathbb E}[\mathbf{x} \, y] - \sum_{j=1}^\infty {\mathbb E}[{\bf x} \, \psi_j({\bf x})] \, {\mathbb E}[y \, \psi_j({\bf x})]$. We then conclude that $\bm{\beta}_0=(\bm{\Sigma}_x)^{-1}\bm{\Sigma}_{xy}$ is invariant to the choice of the orthonormal basis $\{\psi_j(\cdot): j\geq 1\}$.
proof[Proof of Theorem 3.1] (i) It follows that \begin{align*} \widehat{\bm{\beta}}_{\rm SP}=&(\mathbf{X}^\top \mathbf{M}_v \, \mathbf{X})^{-1} \mathbf{X}^\top \mathbf{M}_v \, \mathbf{y}=\bm{\beta}_0 +(\mathbf{X}^\top \mathbf{M}_v \, \mathbf{X})^{-1} \mathbf{X}^\top \mathbf{M}_v \, {\bm\delta} +(\mathbf{X}^\top \mathbf{M}_v \, \mathbf{X})^{-1} \mathbf{X}^\top \mathbf{M}_v \, \mathbf{e}. \end{align*} We then have ${\mathbb E}[\widehat{\bm{\beta}}_{\rm SP}|\mathbf{X}] = \bm{\beta}_0 +(\mathbf{X}^\top \mathbf{M}_v \, \mathbf{X})^{-1} \mathbf{X}^\top \mathbf{M}_v \, \bm{\delta}$. Notice that, as $n\to\infty$, \begin{align*} (\mathbf{X}^\top \mathbf{M}_v \mathbf{X})^{-1} \mathbf{X}^\top \mathbf{M}_v \bm{\delta}=\left(\frac{1}{n} \mathbf{X}^\top \mathbf{M}_v \, \mathbf{X}\right)^{-1} \frac{1}{n} \mathbf{X}^\top \mathbf{M}_v \bm{\delta}\to_P 0, \end{align*} which follows from the proof of part (ii) of Theorem 3.1. The assertion holds. (ii) Notice that $\widehat{\bm{\beta}}_{\rm SP}=\bm{\beta}_0 +(\mathbf{X}^\top \mathbf{M}_v \mathbf{X})^{-1} \mathbf{X}^\top \mathbf{M}_v (\bm{\delta}+\mathbf{e})$ and hence \begin{align*} {\sqrt{n}\left(\frac{\displaystyle1}{\Dn} \mathbf{X}^\top \mathbf{M}_v \mathbf{X}\right)(\widehat{\bm{\beta}}_{\rm SP}-\bm{\beta}_0)= \frac{\displaystyle1}{\sqrt{n}} \mathbf{X}^\top \mathbf{M}_v \, \mathbf{e}+ \frac{\displaystyle1}{\sqrt{n}} \mathbf{X}^\top \mathbf{M}_v \bm{\delta}.} \end{align*} We first consider the convergence of $\frac{\displaystyle1}{\Dn} \mathbf{X}^\top \mathbf{M}_v \mathbf{X}$. Note that \begin{align*} \frac{\displaystyle1}{\Dn} \mathbf{X}^\top \mathbf{M}_v \mathbf{X}=&\frac{\displaystyle1}{\Dn} \mathbf{X}^\top \mathbf{X}-\frac{\displaystyle1}{\Dn} \mathbf{X}^\top \mathbf{V}(\mathbf{V}^\top \mathbf{V})^{-1} \mathbf{V}^\top \mathbf{X}\\ =&\frac{\displaystyle1}{\Dn}\sum_{i=1}^n\mathbf{x}_i\mathbf{x}_i^\top- \frac{\displaystyle1}{\Dn}\sum_{i=1}^n\mathbf{x}_i\mathbf{V}_k(\mathbf{x}_i)^\top \left(\frac{\displaystyle1}{\Dn}\mathbf{V}^\top \mathbf{V}\right)^{-1}\frac{\displaystyle1}{\Dn}\sum_{i=1}^n\mathbf{V}_k(\mathbf{x}_i) \mathbf{x}_i^\top\\ =&{\mathbb E}[\mathbf{x}\mathbf{x}^\top]-{\mathbb E}[\mathbf{x}\mathbf{V}_k(\mathbf{x})^\top] {\mathbb E}[\mathbf{V}_k(\mathbf{x})\mathbf{x}^\top]+o_P(1)\\ \to_P&{\mathbb E}[\mathbf{x}\mathbf{x}^\top]-\sum_{j=1}^\infty{\mathbb E}[\mathbf{x}\psi_j(\mathbf{x})] {\mathbb E}[\psi_j(\mathbf{x})\mathbf{x}^\top]={\bm\Sigma}_x, \end{align*} by Lemma 2.1 as $(k, n)\to(\infty, \infty)$. The asymptotic normality will be derived from $\frac{\displaystyle1}{\displaystyle\sqrt{n}} \mathbf{X}^\top \mathbf{M}_v \, \mathbf{e}$. Because its conditional covariance matrix is $\frac{\displaystyle1}{\Dn}{\mathbb E}[ \mathbf{X}^\top \mathbf{M}_v \, \mathbf{e} \, \mathbf{e}^\top \mathbf{M}_v \mathbf{X}|\mathbf{X}|{\bf X}]=\frac{\displaystyle1}{\Dn} \mathbf{X}^\top \mathbf{M}_v \, \mathbf{\Omega} \, \mathbf{M}_v \mathbf{X}$, and due to the independence data structure, we have \begin{align*} \frac{\displaystyle1}{\sqrt{n}}\left(\frac{\displaystyle1}{\Dn} \mathbf{X}^\top \mathbf{M}_v \, \mathbf{\Omega} \, \mathbf{M}_v \mathbf{X}\right)^{-1/2} \mathbf{X}^\top \mathbf{M}_v \, \mathbf{e}\to_{\mathcal D} N\left(0, {\bf I}_d\right) \end{align*} as $(k, n)\to (\infty, \infty)$. {Meanwhile, we have \begin{align*} &\frac{\displaystyle1}{\Dn} \mathbf{X}^\top \mathbf{M}_v \, \mathbf{\Omega} \, \mathbf{M}_v \mathbf{X} = \frac{\displaystyle1}{\Dn} \mathbf{X}^\top ({\bf I}-\mathbf{P}_v) \, \mathbf{\Omega} \, ({\bf I}-\mathbf{P}_v) \mathbf{X}\\ =&\frac{\displaystyle1}{\Dn} \mathbf{X}^\top \mathbf{\Omega} \mathbf{X}+\frac{\displaystyle1}{\Dn} \mathbf{X}^\top \mathbf{P}_v \, \mathbf{\Omega} \, \mathbf{P}_v\mathbf{X}-\frac{\displaystyle1}{\Dn} \mathbf{X}^\top \mathbf{\Omega} \, \mathbf{P}_v\mathbf{X} -\frac{\displaystyle1}{\Dn} \mathbf{X}^\top \mathbf{P}_v \, \mathbf{\Omega} \, \mathbf{X}\\ =&\frac{\displaystyle1}{\Dn} \mathbf{X}^\top \mathbf{\Omega} \mathbf{X}+\frac{\displaystyle1}{\Dn^3} \mathbf{X}^\top \mathbf{V}\mathbf{V}^\top \, \mathbf{\Omega} \, \mathbf{V}\mathbf{V}^\top\mathbf{X}\\ &-\frac{\displaystyle1}{\Dn^2} \mathbf{X}^\top \mathbf{\Omega} \, \mathbf{V}\mathbf{V}^\top\mathbf{X} -\frac{\displaystyle1}{\Dn^2} \mathbf{X}^\top \mathbf{V}\mathbf{V}^\top \, \mathbf{\Omega} \, \mathbf{X}+o_P(1)\\ =&\frac{\displaystyle1}{\Dn}\sum_{i=1}^n \mathbf{x}_i\mathbf{x}^\top_i \sigma_i^2+ \frac{\displaystyle1}{\Dn}\sum_{i=1}^n \mathbf{x}_i\mathbf{V}_k(\mathbf{x}_i)^\top \frac{\displaystyle1}{\Dn}\sum_{i=1}^n \mathbf{V}_k(\mathbf{x}_i)\mathbf{V}_k(\mathbf{x}_i)^\top \sigma_i^2 \frac{\displaystyle1}{\Dn}\sum_{i=1}^n \mathbf{V}_k(\mathbf{x}_i)\mathbf{x}^\top_i \\ &-\frac{\displaystyle1}{\Dn}\sum_{i=1}^n \mathbf{x}_i\mathbf{V}_k(\mathbf{x}_i)^\top \sigma_i^2\frac{\displaystyle1}{\Dn}\sum_{i=1}^n \mathbf{V}_k(\mathbf{x}_i)\mathbf{x}_i^\top-\frac{\displaystyle1}{\Dn}\sum_{i=1}^n \mathbf{x}_i\mathbf{V}_k(\mathbf{x}_i)^\top \frac{\displaystyle1}{\Dn}\sum_{i=1}^n \mathbf{V}_k(\mathbf{x}_i)\mathbf{x}_i^\top \sigma_i^2\\ =&\left({\mathbb E}[\mathbf{x}\mathbf{x}^\top]-{\mathbb E}[\mathbf{x}\mathbf{V}_k(\mathbf{x})^\top]{\mathbb E}[\mathbf{V}_k(\mathbf{x})\mathbf{x}^\top]\right)\frac{\displaystyle1}{\Dn}\sum_{i=1}^n\sigma_i^2+o_P(1)\to_P\bm{\Sigma}_x\bar{\sigma}^2, \end{align*} as $n,k\to\infty$, where we have used the following approximations: \begin{align*} &\frac{\displaystyle1}{\Dn}\sum_{i=1}^n \mathbf{x}_i\mathbf{x}^\top_i \sigma_i^2={\mathbb E}[\mathbf{x}\mathbf{x}^\top]\frac{\displaystyle1}{\Dn}\sum_{i=1}^n\sigma_i^2+O_P(n^{-1/2}),\\ &\frac{\displaystyle1}{\Dn}\sum_{i=1}^n \mathbf{x}_i\mathbf{V}_k(\mathbf{x}_i)^\top ={\mathbb E}[\mathbf{x}\mathbf{V}_k(\mathbf{x})^\top] +O_P(k^{1/2}n^{-1/2}),\\ &\frac{\displaystyle1}{\Dn}\sum_{i=1}^n \mathbf{V}_k(\mathbf{x}_i)\mathbf{V}_k(\mathbf{x}_i)^\top \sigma_i^2=I_k\frac{\displaystyle1}{\Dn}\sum_{i=1}^n \sigma_i^2+O(kn^{-1/2}),\\ &\frac{\displaystyle1}{\Dn}\sum_{i=1}^n \mathbf{V}_k(\mathbf{x}_i)\mathbf{x}_i^\top \sigma_i^2={\mathbb E}[\mathbf{V}_k(\mathbf{x})\mathbf{x}^\top] \frac{\displaystyle1}{\Dn}\sum_{i=1}^n \sigma_i^2+O_P(k^{1/2}n^{-1/2}), \end{align*} provided that $\sum_{i=1}^n \sigma_i^4=O(n)$. Since they are quite easily verified, we omit their proofs. } It remains to show $\frac{\displaystyle1}{\displaystyle\sqrt{n}} \mathbf{X}^\top \mathbf{M}_v \bm{\delta}\to_P0$. Because $\mathbf{M}_v$ is idempotent, its eigenvalues are either one or zero. It follows that \begin{eqnarray} && \frac{1}{\sqrt{n}}\|\mathbf{X}^\top \mathbf{M}_v \bm{\delta}\|\leq \frac{\displaystyle1}{\sqrt{n}}\|\mathbf{X}\|\|\bm{\delta}\| =\sqrt{\frac{1}{n}\sum_{i=1}^n\|\mathbf{x}_i\|^2} \sqrt{\sum_{i=1}^n |\delta_k(\mathbf{x}_i)|^2}, \end{eqnarray} and ${\mathbb E}[\delta_k^2(\mathbf{\bf x}_i)]=\int_{{\mathcal X}} \delta_k^2(\mathbf{x}) f(\mathbf{x})d \mathbf{x}=o(k^{-2s/d})$ by Lemma A.5 in dlp2021. This implies $\frac{\displaystyle1}{\displaystyle\sqrt{n}} \mathbf{X}^\top \mathbf{M}_v \bm{\delta}=O_P(\sqrt{n \, k^{-2s/d}})=o_P(1)$ by Assumption 3.1. (iii)\ Notice that \begin{align*} &\frac{\displaystyle1}{\Dn} \mathbf{X}^\top \mathbf{M}_v \,(\widehat{\bf \Omega}- \mathbf{\Omega}) \, \mathbf{M}_v \mathbf{X} = \frac{\displaystyle1}{\Dn} \mathbf{X}^\top ({\bf I}-\mathbf{P}_v) \, (\widehat{\bf \Omega}- \mathbf{\Omega}) \, ({\bf I}-\mathbf{P}_v) \mathbf{X}\\ =&\frac{\displaystyle1}{\Dn} \mathbf{X}^\top (\widehat{\bf \Omega}- \mathbf{\Omega}) \mathbf{X}+\frac{\displaystyle1}{\Dn} \mathbf{X}^\top \mathbf{P}_v (\widehat{\bf \Omega}- \mathbf{\Omega}) \mathbf{P}_v\mathbf{X}-\frac{\displaystyle1}{\Dn} \mathbf{X}^\top (\widehat{\bf \Omega}- \mathbf{\Omega}) \mathbf{P}_v\mathbf{X} -\frac{\displaystyle1}{\Dn} \mathbf{X}^\top \mathbf{P}_v (\widehat{\bf \Omega}- \mathbf{\Omega}) \mathbf{X}\\ =&\frac{\displaystyle1}{\Dn} \mathbf{X}^\top (\widehat{\bf \Omega}- \mathbf{\Omega})\mathbf{X}+\frac{\displaystyle1}{\Dn^3} \mathbf{X}^\top \mathbf{V}\mathbf{V}^\top (\widehat{\bf \Omega}- \mathbf{\Omega}) \mathbf{V}\mathbf{V}^\top\mathbf{X}\\ &-\frac{\displaystyle1}{\Dn^2} \mathbf{X}^\top (\widehat{\bf \Omega}- \mathbf{\Omega}) \mathbf{V}\mathbf{V}^\top\mathbf{X} -\frac{\displaystyle1}{\Dn^2} \mathbf{X}^\top \mathbf{V}\mathbf{V}^\top (\widehat{\bf \Omega}- \mathbf{\Omega}) \mathbf{X}+o_P(1)\\ =&\frac{\displaystyle1}{\Dn}\sum_{i=1}^n \mathbf{x}_i\mathbf{x}^\top_i(\widehat{e}_i^2-\sigma_i^2)+ \frac{\displaystyle1}{\Dn}\sum_{i=1}^n \mathbf{x}_i\mathbf{V}_k(\mathbf{x}_i)^\top \frac{\displaystyle1}{\Dn}\sum_{i=1}^n \mathbf{V}_k(\mathbf{x}_i)\mathbf{V}_k(\mathbf{x}_i)^\top (\widehat{e}_i^2- \sigma_i^2) \frac{\displaystyle1}{\Dn}\sum_{i=1}^n \mathbf{V}_k(\mathbf{x}_i)\mathbf{x}^\top_i \\ &-\frac{\displaystyle1}{\Dn}\sum_{i=1}^n \mathbf{x}_i\mathbf{V}_k(\mathbf{x}_i)^\top(\widehat{e}_i^2-\sigma_i^2)\frac{\displaystyle1}{\Dn}\sum_{i=1}^n \mathbf{V}_k(\mathbf{x}_i)\mathbf{x}_i^\top-\frac{\displaystyle1}{\Dn}\sum_{i=1}^n \mathbf{x}_i\mathbf{V}_k(\mathbf{x}_i)^\top \frac{\displaystyle1}{\Dn}\sum_{i=1}^n \mathbf{V}_k(\mathbf{x}_i)\mathbf{x}_i^\top (\widehat{e}_i^2-\sigma_i^2). \end{align*} It suffices to show \begin{align} I_1\equiv&\frac{\displaystyle1}{\Dn}\sum_{i=1}^n \mathbf{x}_i\mathbf{x}^\top_i(\widehat{e}_i^2-\sigma_i^2)=o_P(1),\\ I_2\equiv&\frac{\displaystyle1}{\Dn}\sum_{i=1}^n \mathbf{V}_k(\mathbf{x}_i)\mathbf{V}_k(\mathbf{x}_i)^\top (\widehat{e}_i^2-\sigma_i^2)=o_P(1),\\ I_3\equiv&\frac{\displaystyle1}{\Dn}\sum_{i=1}^n \mathbf{V}_k(\mathbf{x}_i)\mathbf{x}_i^\top (\widehat{e}_i^2-\sigma_i^2)=o_P(1). \end{align} To do this, note that \begin{eqnarray} \widehat{e}_i^2 &=&(y_i-{\bf x}_i^\top \widehat{\bm{\beta}}_{\rm SP}-\mathbf{V}_k({\bf x}_i)^\top \widehat{\bm{\gamma}})^2 \notag \\ &=& [e_i -{\bf x}_i^\top( \widehat{\bm{\beta}}_{\rm SP}-{\bm \beta}_0)-\mathbf{V}_k({\bf x}_i)^\top ({\bm \gamma}-\widehat{\bm{\gamma}})+\delta_k({\bf x}_i)]^2 \notag \\ &=& e_i^2 -2e_i [{\bf x}_i^\top( \widehat{\bm{\beta}}_{\rm SP}-{\bm \beta}_0)+\mathbf{V}_k({\bf x}_i)^\top ({\bm \gamma}-\widehat{\bm{\gamma}})-\delta_k({\bf x}_i)] \notag\\ &&+[{\bf x}_i^\top( \widehat{\bm{\beta}}_{\rm SP}-{\bm \beta}_0)+\mathbf{V}_k({\bf x}_i)^\top ({\bm \gamma}- \widehat{\bm{\gamma}}) -\delta_k({\bf x}_i)]^2. \nonumber \end{eqnarray} Thus, for (ref), \begin{eqnarray} && I_1=\frac{\displaystyle1}{\Dn}\sum_{i=1}^n \mathbf{x}_i\mathbf{x}^\top_i(\widehat{e}_i^2-\sigma_i^2)\\ && =\frac{\displaystyle1}{\Dn}\sum_{i=1}^n \mathbf{x}_i\mathbf{x}^\top_i(e_i^2-\sigma_i^2)+\frac{\displaystyle2}{\Dn}\sum_{i=1}^n \mathbf{x}_i\mathbf{x}^\top_ie_i[{\bf x}_i^\top( \widehat{\bm{\beta}}_{\rm SP}-{\bm \beta}_0)-\mathbf{V}_k({\bf x}_i)^\top ({\bm \gamma}-\widehat{\bm{\gamma}})+\delta_k({\bf x}_i)]\nonumber\\ &&+\frac{\displaystyle1}{\Dn}\sum_{i=1}^n \mathbf{x}_i\mathbf{x}^\top_i[{\bf x}_i^\top( \widehat{\bm{\beta}}_{\rm SP}-{\bm \beta}_0)-\mathbf{V}_k({\bf x}_i)^\top (\widehat{\bm{\gamma}}-{\bm \gamma})+\delta_k({\bf x}_i)]^2\equiv I_{11}+I_{12}+I_{13}. \nonumber \end{eqnarray} Moreover, \begin{align*} {\mathbb E}[\|I_{11}\|^2]=&\frac{\displaystyle1}{\Dn^2}\sum_{i=1}^n {\mathbb E}[\|\mathbf{x}_i\mathbf{x}^\top_i\|^2(e_i^2-\sigma_i^2)^2]\\ =&\frac{\displaystyle1}{\Dn^2}\sum_{i=1}^n {\mathbb E}\{\|\mathbf{x}_i\mathbf{x}^\top_i\|^2{\mathbb E}[(e_i^2-\sigma_i^2)^2|{\bf x}_i]\} \le\frac{\displaystyle1}{\Dn}\mu_4{\mathbb E}[\|{\bf x}\|^4]=o(1), \end{align*} where $\mu_4={\mathbb E}[e_i^4|{\bf x}_i]$; and \begin{eqnarray} &&\|I_{13}\|\le \frac{\displaystyle1}{\Dn}\sum_{i=1}^n \|\mathbf{x}_i\|^2[{\bf x}_i^\top( \widehat{\bm{\beta}}_{\rm SP}-{\bm \beta}_0)+\mathbf{V}_k({\bf x}_i)^\top ({\bm \gamma}-\widehat{\bm{\gamma}})-\delta_k({\bf x}_i)]^2 \nonumber\\ &&\le \frac{\displaystyle6}{\Dn}\sum_{i=1}^n \|\mathbf{x}_i\|^2[{\bf x}_i^\top( \widehat{\bm{\beta}}_{\rm SP}-{\bm \beta}_0)]^2+\frac{\displaystyle6}{\Dn}\sum_{i=1}^n \|\mathbf{x}_i\|^2[\mathbf{V}_k({\bf x}_i)^\top (\widehat{\bm{\gamma}}-{\bm \gamma})]^2 \nonumber\\ &&+\frac{\displaystyle6}{\Dn}\sum_{i=1}^n \|\mathbf{x}_i\|^2\delta_k^2({\bf x}_i)\le \| \widehat{\bm{\beta}}_{\rm SP}-{\bm \beta}_0)\|^2\frac{\displaystyle6}{\Dn}\sum_{i=1}^n \|\mathbf{x}_i\|^4+\|\widehat{\bm{\gamma}}-{\bm \gamma}\|^2\frac{\displaystyle6}{\Dn}\sum_{i=1}^n \|\mathbf{x}_i\|^2\|\mathbf{V}_k({\bf x}_i)\|^2 \nonumber\\ &&+\frac{\displaystyle6}{\Dn}\sum_{i=1}^n \|\mathbf{x}_i\|^2\delta_k^2({\bf x}_i) = O_P(\|\widehat{\bm{\beta}}_{\rm SP}-{\bm \beta}_0)\|^2)+O_P(\|\widehat{\bm{\gamma}}-{\bm \gamma}\|^2) {\mathbb E}[\|\mathbf{x}_i\|^2\|\mathbf{V}_k({\bf x}_i)\|^2] \nonumber\\ && +O_P({\mathbb E}[\|\mathbf{x}_i\|^2\delta_k^2({\bf x}_i)]) = O_P(\|\widehat{\bm{\beta}}_{\rm SP}-{\bm \beta}_0)\|^2)+O_P(k\|\widehat{\bm{\gamma}}-{\bm \gamma}\|^2)+ O_P\left(\{{\mathbb E}[\delta_k^4({\bf x}_i)]\}^{1/2}\right) =o_P(1),\nonumber \end{eqnarray} because $\sqrt{k}(\widehat{\bm{\gamma}}-{\bm \gamma})=\sqrt{k}\frac{\displaystyle1}{\Dn}{\bf V}^\top {\bf X}(\widehat{\bm{\beta}}_{\rm SP}-{\bm \beta}_0)+\sqrt{k}\frac{\displaystyle1}{\Dn}{\bf V}^\top{\bf e}+\sqrt{k}\frac{\displaystyle1}{\Dn}{\bf V}^\top{\bm \delta}=O_P(kn^{-1/2})+O_P(kk^{-s/d})=o_P(1)$ due to Assumption 3.1, and as $k\rightarrow \infty$ \begin{eqnarray} &&{\mathbb E}[\delta_k^4(\mathbf{x}_i)]=\int \delta_k^4(x)f_{\bf x}(x)dx= \sum_{j=k+1}^\infty \gamma_j^4\int \psi_j^4(x) f_{\bf x}(x)dx \nonumber\\ &&+6 \sum_{s,j=k+1}^\infty \gamma_j^2\gamma_s^2\int \psi_j^2(x)\psi_s^2(x) f_{\bf x}(x)dx\le C_1\sum_{j=k+1}^\infty \gamma_j^4+C_2\sum_{s,j=k+1}^\infty \gamma_j^2\gamma_s^2=o(1).\nonumber \end{eqnarray} The derivation of $I_{12}=o_P(1)$ follows similarly, so we omit the detail. For (ref), \begin{eqnarray} && I_2=\frac{\displaystyle1}{\Dn}\sum_{i=1}^n \mathbf{V}_k(\mathbf{x}_i)\mathbf{V}_k(\mathbf{x}_i)^\top (\widehat{e}_i^2-\sigma_i^2)=\frac{\displaystyle1}{\Dn}\sum_{i=1}^n \mathbf{V}_k(\mathbf{x}_i)\mathbf{V}_k(\mathbf{x}_i)^\top (e_i^2-\sigma_i^2) \nonumber\\ &&-\frac{\displaystyle2}{\Dn}\sum_{i=1}^n \mathbf{V}_k(\mathbf{x}_i)\mathbf{V}_k(\mathbf{x}_i)^\top e_i [{\bf x}_i^\top( \widehat{\bm{\beta}}_{\rm SP}-{\bm \beta}_0)+\mathbf{V}_k({\bf x}_i)^\top ({\bm \gamma}-\widehat{\bm{\gamma}})-\delta_k({\bf x}_i)] \nonumber\\ &&+\frac{\displaystyle1}{\Dn}\sum_{i=1}^n \mathbf{V}_k(\mathbf{x}_i)\mathbf{V}_k(\mathbf{x}_i)^\top[{\bf x}_i^\top( \widehat{\bm{\beta}}_{\rm SP}-{\bm \beta}_0)+\mathbf{V}_k({\bf x}_i)^\top ({\bm \gamma}- \widehat{\bm{\gamma}}) -\delta_k({\bf x}_i)]^2 =I_{21}+I_{22}+I_{23}. \nonumber \end{eqnarray} Note that \begin{align*} {\mathbb E}[\|I_{21}\|^2]=&\frac{\displaystyle1}{\Dn^2}\sum_{i=1}^n {\mathbb E}\|\mathbf{V}_k(\mathbf{x}_i)\mathbf{V}_k(\mathbf{x}_i)^\top\|^2 (e_i^2-\sigma_i^2)^2 \le \frac{\mu_4}{\Dn}{\mathbb E}\|\mathbf{V}_k(\mathbf{x})\|^4=C\frac{\displaystyle1}{\Dn}k^2=o(1); \nonumber\\ \|I_{23}\|\le&\frac{\displaystyle1}{\Dn}\sum_{i=1}^n \|\mathbf{V}_k(\mathbf{x}_i)\|^2[{\bf x}_i^\top( \widehat{\bm{\beta}}_{\rm SP}-{\bm \beta}_0)]^2+\frac{\displaystyle1}{\Dn}\sum_{i=1}^n \|\mathbf{V}_k(\mathbf{x}_i)\|^2\delta_k^2({\bf x}_i)\\ &+\frac{\displaystyle1}{\Dn}\sum_{i=1}^n \|\mathbf{V}_k(\mathbf{x}_i)\|^2[\mathbf{V}_k({\bf x}_i)^\top ({\bm \gamma}- \widehat{\bm{\gamma}})]^2\\ \le&\|\widehat{\bm{\beta}}_{\rm SP}-{\bm \beta}_0\|^2O_P({\mathbb E} \|\mathbf{V}_k(\mathbf{x})\|^2\|{\bf x}\|^2)+O_P({\mathbb E}\|\mathbf{V}_k(\mathbf{x})\|^2\delta_k^2({\bf x}))\\ &+\|{\bm \gamma}- \widehat{\bm{\gamma}}\|^2O_P({\mathbb E} \|\mathbf{V}_k(\mathbf{x})\|^4)\\ =&O_P(kn^{-1})+O_P(k^2k^{-2s/d})+O_P(k^{3/2}n^{-1/2})=o_P(1), \end{align*} due to Assumption 3.1 again. We also omit the verification of $I_{22}$ due to the same reason. Meanwhile, as the condition ensuring $I_3$ to hold is no more restrictive than that for $I_{22}$, $I_3$ holds automatically. We therefore have completed the proof of Theorem 3.1. Actually, it should be noted that the proof of Theorem 3.1 follows trivially from those of Theorems A.1 and A.2 in the case of homoskedasticity ${\mathbb E}[e_i^2|\mathbf{x}_i] =\sigma_e^2$ (a.s.).
proof[Proofs of Theorems 3.2 and 3.3] The proofs of Theorems 3.2 and 3.3 follow respectively from those of Theorems A.1 and A.2 of the online supplementary document, in which we assume $\sigma_i^2 = {\mathbb E}[e_i^2|\mathbf{x}_i] =\sigma_e^2$ (a.s.) for notational simplicity. As may be seen from the derivations in Appendix A.4.3 and then the proof of Theorem 3.1 above, we need only to make some minor changes to those involved in the relevant places in the derivations of the proofs of Theorems 3.2 and 3.3. We therefore omit the repetitions.

{

proof[Proof of Theorem 5.1] (i)-(ii). Before proceeding, we justify some basic facts. Let $\pmb{\Omega}_k =\begin{pmatrix} \mathbf{A}_{11} & \mathbf{A}_{12} \\ \mathbf{A}_{12}^\top & \mathbf{I}_k \\ \end{pmatrix}$, where $\mathbf{A}_{11}={\mathbb E}[\mathbf{x}\mathbf{x}^\top]$, $\mathbf{A}_{12}\coloneqq (\pmb{\eta}_1,\cdots , \pmb{\eta}_{k})$, and $\pmb{\eta}_\ell={\mathbb E}[\mathbf{x}\phi_\ell(\mathbf{x})]$. To show the invertibility of $\pmb{\Omega}_k$, we note that $ \pmb{\Omega}_k =\begin{pmatrix} \mathbf{A}_{11} &\mathbf{0} \\ \mathbf{0} & \mathbf{I}_k \\ \end{pmatrix}+\begin{pmatrix} \mathbf{0} & \mathbf{A}_{12} \\ \mathbf{A}_{12}^\top & \mathbf{0} \\ \end{pmatrix}$, so by Weyl's inequality, \begin{eqnarray} \rho_{\min}(\pmb{\Omega}_k) & \ge & \rho_{\min}(\mathbf{A}_{11})\wedge 1 +\rho_{\min}^{1/2}\left(\begin{pmatrix} \mathbf{0} & \mathbf{A}_{12} \\ \mathbf{A}_{12}^\top & \mathbf{0} \\ \end{pmatrix}\begin{pmatrix} \mathbf{0} & \mathbf{A}_{12} \\ \mathbf{A}_{12}^\top & \mathbf{0} \\ \end{pmatrix}\right)\notag \\ &=&\rho_{\min}(\mathbf{A}_{11})\wedge 1 +\rho_{\min}^{1/2}\left(\begin{pmatrix} \mathbf{A}_{12}\mathbf{A}_{12}^\top & \mathbf{0} \\ \mathbf{0} & \mathbf{A}_{12}^\top\mathbf{A}_{12} \\ \end{pmatrix} \right) >0, \nonumber \end{eqnarray} where $\rho_{\min}(\cdot)$ stands for the minimum eigenvalue value. Therefore, for $\forall(\boldsymbol{\beta}, \boldsymbol{\xi}_k)\in \mathbb{A}$ and $\operatorname*{\normalfont\textrm{vec}}(\boldsymbol{\beta}, \boldsymbol{\xi}_k)\ne \mathbf{0}$, we have $\operatorname*{\normalfont\textrm{vec}}(\boldsymbol{\beta}, \boldsymbol{\xi}_k)^\top \pmb{\Omega}_k\operatorname*{\normalfont\textrm{vec}}(\boldsymbol{\beta}, \boldsymbol{\xi}_k)\ge \alpha \|\operatorname*{\normalfont\textrm{vec}}(\boldsymbol{\beta}, \boldsymbol{\xi}_k)\|^2$ uniformly in $k$, where $\mathbb{A} \coloneqq \{ \operatorname*{\normalfont\textrm{vec}}(\boldsymbol{\beta}, \boldsymbol{\xi}_k)\mid \|\operatorname*{\normalfont\textrm{vec}}(\boldsymbol{\beta}, \mathbf{S}_{\bar{\mathscr{C}}} \boldsymbol{\xi}_k ) \|_1\le 3 \|\operatorname*{\normalfont\textrm{vec}}(\boldsymbol{\beta}, \mathbf{S}_{\mathscr{C}} \boldsymbol{\xi}_k )\|_1\}$, and $\alpha>0$ is a fixed positive constant, in which $\mathbf{S}_{\mathscr{C}}$ and $\mathbf{S}_{\bar{\mathscr{C}}} $ are respectively $k_0\times k$ and $(k-k_0)\times k$ selection matrices selecting the elements indexed by $\mathscr{C}$ and $\bar{\mathscr{C}}$. In this case the so-called restricted eigenvalue condition is automatically fulfilled, which has been fully discussed in the literature, such as BRT2009 and raskutti10a. We now proceed. To avoid introducing any new notation, in a similar way to (3.2) of Section 3, we consider a vector form of model (5.1): \begin{eqnarray} \mathbf{y} = \mathbf{X}\boldsymbol{\beta}_0 + \mathbf{V}\boldsymbol{\xi} +\boldsymbol{\delta} +\mathbf{e}. \end{eqnarray} We then define the following function: $Q(\boldsymbol{\beta}, \boldsymbol{\xi}_k) =\frac{1}{n}\|\mathbf{y}- \mathbf{X}\boldsymbol{\beta} - \mathbf{V}\boldsymbol{\xi}_k\|^2 + \lambda \|\boldsymbol{\xi}_k\|_1$ to conduct the LASSO regression, where $\|\cdot \|_1$ stands for $L_1$ norm, $\boldsymbol{\beta}$ and $\boldsymbol{\xi}_k$ are generic vectors with dimensions being the same as $\boldsymbol{\beta}_0$ and $\boldsymbol{\xi}$ respectively, in which $\boldsymbol{\xi}$ collects all $\xi_\ell$'s stated in (5.2). The LASSO estimates of $\boldsymbol{\beta}_0$ and $\boldsymbol{\xi}$ are given by $(\widehat{\boldsymbol{\beta}}^*, \widehat{\boldsymbol{\xi}}^*) =\operatorname*{\arg\!\min} Q(\boldsymbol{\beta}, \boldsymbol{\xi}_k)$. By definition, we have \begin{eqnarray} Q(\widehat{\boldsymbol{\beta}}^*, \widehat{\boldsymbol{\xi}}^*)\le Q(\boldsymbol{\beta}_0, \boldsymbol{\xi}) = \frac{1}{n}\|\mathbf{e} + \boldsymbol{\delta}\| + \lambda \| \boldsymbol{\xi} \|_1. \end{eqnarray} To further expand the left hand side of (ref), we study the following term. Note that \begin{eqnarray*} \frac{1}{n}\left| (\mathbf{e} + \boldsymbol{\delta})^\top [\mathbf{X}, \mathbf{V}] \operatorname*{\normalfontvec}(\widehat{\boldsymbol{\beta}}^*- \boldsymbol{\beta}_0, \widehat{\boldsymbol{\xi}}^* - \boldsymbol{\xi} )\right| \le \frac{1}{n} \| (\mathbf{e} + \boldsymbol{\delta})^\top [\mathbf{X}, \mathbf{V}] \|_\infty \|\operatorname*{\normalfontvec}(\widehat{\boldsymbol{\beta}}^* - \boldsymbol{\beta}_0,\widehat{\boldsymbol{\xi}}^* - \boldsymbol{\xi}) \|_1, \end{eqnarray*} where $\|\cdot \|_\infty$ stands for the maximum norm of a vector. We focus on $\frac{1}{n} \| (\mathbf{e} + \boldsymbol{\delta})^\top [\mathbf{X}, \mathbf{V}] \|_\infty$ in what follows. Note that $\boldsymbol{\phi}_\ell \coloneqq (\phi_\ell(\mathbf{x}_1),\ldots, \phi_\ell(\mathbf{x}_n))^\top$ is the $\ell^{th}$ column of $\mathbf{V}$ and then derive \begin{eqnarray*} \max_{\ell}\frac{1}{n}|\boldsymbol{\delta}^\top \boldsymbol{\phi}_\ell|\le \max_{\ell}\frac{1}{n}\sum_{i=1}^n |\delta_k(\mathbf{x}_i) \phi_\ell(\mathbf{x}_i)| \le O(1)\frac{1}{n}\sum_{i=1}^n |\delta_k(\mathbf{x}_i) |=o_P(k^{-s/d}), \end{eqnarray*} where the second inequality follows from $\phi_\ell(\cdot)$ being uniformly bounded, and the last step follows form the fact that ${\mathbb E}|\delta_k(\mathbf{x}_i) |\le \{{\mathbb E}|\delta_k(\mathbf{x}_i) |^2\}^{1/2} =o(k^{-s/d})$ in which the first inequality follows from the moment inequality, and the last step has been proved in Theorem 3.1. Similarly, we obtain $ \frac{1}{n}\|\boldsymbol{\delta}^\top \mathbf{X}\| =o_P(k^{-s/d})$. Thus, we can conclude \begin{eqnarray} \frac{1}{n} \|\boldsymbol{\delta}^\top [\mathbf{X}, \mathbf{V}]\|_\infty=o_P(k^{-s/d}). \end{eqnarray} We consider $\frac{1}{n} \mathbf{e}^\top \boldsymbol{\phi}_\ell =\frac{1}{n}\sum_{i=1}^n \phi_\ell (\mathbf{x}_i) \, e_i$. In order to do so, we let $\epsilon = c_0\sqrt{n\log k}$ and write \begin{eqnarray*} &&{\rm P}\left(\sum_{i=1}^n \phi_\ell (\mathbf{x}_i)e_i\ge \epsilon\right) \le c_1 \, \sum_{i=1}^n\|e_i\|_J^J \ \epsilon^{-J} + 2\exp\left(- c_2\epsilon^2 \ \left(\sum_{i=1}^n\| e_i\|_2^2\right)^{-1}\right) \notag \\ &=& c_1 \frac{n}{\epsilon^J} + 2\exp\left( \frac{- c_2\epsilon^2}{n}\right) = \frac{c_1}{n^{J-1}(\log k)^{J/2}}+2\, e^{- c_2\log k}, \end{eqnarray*} where the first inequality follows from Corollary 1.8 in nagaev1979large, and we may let $J=4$ as in the condition of Theorem 2.1. Additionally, we require $c_2>1$ which is achievable by varying the value of $c_0$. Similarly, we obtain that ${\rm P}\left(\sum_{i=1}^n x_{i,j} e_i\ge \epsilon\right) = \frac{c_1}{n^{J-1}(\log k)^{J/2}}+2\exp\left( - c_2\log k\right)$, where $x_{i,j}$ stands for the $j^{th}$ element of $\mathbf{x}_i$. Thus, under the condition $\frac{k}{n^{J-1}(\log k)^{J/2}}\to 0$, \begin{eqnarray} \frac{1}{n} \| \mathbf{e}^\top [\mathbf{X}, \mathbf{V}] \|_\infty = O_P\left(\frac{\sqrt{\log k}}{\sqrt{n}}\right). \end{eqnarray} By (ref) and (ref), we obtain that \begin{eqnarray} \frac{1}{n} \| (\mathbf{e} + \boldsymbol{\delta})^\top [\mathbf{X}, \mathbf{V}] \|_\infty=O_P\left(\frac{\sqrt{\log k}}{\sqrt{n}} \vee k^{-s/d}\right)=O_P\left(\frac{\sqrt{\log k}}{\sqrt{n}} \right), \end{eqnarray} where the last step is due to Assumption 3.1(iv). Then with probability approaching 1, we have $\frac{1}{n} \| (\mathbf{e} + \boldsymbol{\delta})^\top [\mathbf{X}, \mathbf{V}] \|_\infty\le \frac{\lambda}{4}$, where $\lambda = c^* \frac{\sqrt{\log k}}{\sqrt{n}} $ with $c^*$ being a sufficiently large number. Then, we have $\frac{1}{n}\left| (\mathbf{e} + \boldsymbol{\delta})^\top [\mathbf{X}, \mathbf{V}] \operatorname*{\normalfont\textrm{vec}}(\boldsymbol{\beta}_0 - \widehat{\boldsymbol{\beta}}^*, \boldsymbol{\xi} -\widehat{\boldsymbol{\xi}}^* )\right| \le \frac{\lambda}{4} \|\operatorname*{\normalfont\textrm{vec}}(\boldsymbol{\beta}_0 - \widehat{\boldsymbol{\beta}}^*, \boldsymbol{\xi} -\widehat{\boldsymbol{\xi}}^* )\|_1$. Therefore, (ref) reduces to \begin{eqnarray} \frac{1}{n} \| [\mathbf{X}, \mathbf{V}_k] \operatorname*{\normalfontvec}(\boldsymbol{\beta}_0 - \widehat{\boldsymbol{\beta}}^*, \boldsymbol{\xi} -\widehat{\boldsymbol{\xi}}^* ) \|^2 + \lambda\|\widehat{\boldsymbol{\xi}}^* \|_1 \le \frac{\lambda}{2} \|\operatorname*{\normalfontvec}(\boldsymbol{\beta}_0 - \widehat{\boldsymbol{\beta}}^*, \boldsymbol{\xi} -\widehat{\boldsymbol{\xi}}^* )\|_1 + \lambda \| \widehat{\boldsymbol{\xi}}^* \|_1. \end{eqnarray} Recall the definitions of $\mathbf{S}_{\mathscr{C}}$ and $\mathbf{S}_{\bar{\mathscr{C}}}$ in Assumption 5.1. Simple algebra shows that \begin{eqnarray*} &&\|\operatorname*{\normalfontvec}(\boldsymbol{\beta}_0 - \widehat{\boldsymbol{\beta}}^*, \boldsymbol{\xi} -\widehat{\boldsymbol{\xi}}^* )\|_1 +\|\operatorname*{\normalfontvec}( \boldsymbol{\beta}_0, \boldsymbol{\xi} )\|_1 -\|\operatorname*{\normalfont\textrm{vec}}( \widehat{\boldsymbol{\beta}}^*, \widehat{\boldsymbol{\xi}}^* )\|_1 \le 2\|\operatorname*{\normalfont\textrm{vec}}(\boldsymbol{\beta}_0 - \widehat{\boldsymbol{\beta}}^*, \mathbf{S}_{\mathscr{C}} ( \boldsymbol{\xi} -\widehat{\boldsymbol{\xi}}^*) )\|_1, \end{eqnarray*} so we can rewrite (ref) as $\frac{1}{n} \| [\mathbf{X}, \mathbf{V}] \operatorname*{\normalfont\textrm{vec}}(\boldsymbol{\beta}_0 - \widehat{\boldsymbol{\beta}}^*, \boldsymbol{\xi} -\widehat{\boldsymbol{\xi}}^* ) \|^2 +\frac{\lambda}{2} \|\operatorname*{\normalfont\textrm{vec}}(\boldsymbol{\beta}_0 - \widehat{\boldsymbol{\beta}}^*, \boldsymbol{\xi} -\widehat{\boldsymbol{\xi}}^* )\|_1 \leq 2\lambda\|\operatorname*{\normalfont\textrm{vec}}(\boldsymbol{\beta}_0 - \widehat{\boldsymbol{\beta}}^*, \mathbf{S}_{\mathscr{C}} ( \boldsymbol{\xi} -\widehat{\boldsymbol{\xi}}^*) )\|_1$. It further yields \begin{eqnarray} && \frac{1}{n} \| [\mathbf{X}, \mathbf{V}] \operatorname*{\normalfont\textrm{vec}}(\boldsymbol{\beta}_0 - \widehat{\boldsymbol{\beta}}^*, \boldsymbol{\xi} -\widehat{\boldsymbol{\xi}}^* ) \|^2 +\frac{\lambda}{2} \|\operatorname*{\normalfont\textrm{vec}}(\boldsymbol{\beta}_0 - \widehat{\boldsymbol{\beta}}^*, \mathbf{S}_{\bar{\mathscr{C}}} ( \boldsymbol{\xi} -\widehat{\boldsymbol{\xi}}^*) )\|_1 \notag \\ &\le & \frac{3}{2}\lambda\|\operatorname*{\normalfont\textrm{vec}}(\boldsymbol{\beta}_0 - \widehat{\boldsymbol{\beta}}^*, \mathbf{S}_{\mathscr{C}} ( \boldsymbol{\xi} -\widehat{\boldsymbol{\xi}}^*) )\|_1. \end{eqnarray} Finally, we focus on $\frac{1}{n} \| [\mathbf{X}, \mathbf{V}] \operatorname*{\normalfont\textrm{vec}}(\boldsymbol{\beta}_0 - \widehat{\boldsymbol{\beta}}^*, \boldsymbol{\xi} -\widehat{\boldsymbol{\xi}}^* ) \|^2$ of (ref). As explained in the beginning of the proof, we have for $\forall (\boldsymbol{\beta}, \boldsymbol{\xi}_k)\in \mathbb{A}$ \begin{eqnarray} \|\operatorname*{\normalfont\textrm{vec}}(\boldsymbol{\beta}, \boldsymbol{\xi}_k ) \|_1 &=& \|\operatorname*{\normalfont\textrm{vec}}(\boldsymbol{\beta}, \mathbf{S}_{\mathscr{C}} \boldsymbol{\xi}_k) \|_1 +\|\operatorname*{\normalfont\textrm{vec}}(\boldsymbol{\beta}, \mathbf{S}_{\bar{\mathscr{C}}}\boldsymbol{\xi}_k) \|_1 \notag \\ &\le & 4 \|\operatorname*{\normalfont\textrm{vec}}(\boldsymbol{\beta}, \mathbf{S}_{\mathscr{C}} \boldsymbol{\xi}_k) \|_1\le 4\sqrt{k_0}\| \operatorname*{\normalfont\textrm{vec}}(\boldsymbol{\beta}, \mathbf{S}_{\mathscr{C}} \boldsymbol{\xi}_k) \| . \end{eqnarray} Thus, with probability approaching 1, \begin{eqnarray} &&\frac{1}{n} \| [\mathbf{X}, \mathbf{V}] \operatorname*{\normalfont\textrm{vec}}(\boldsymbol{\beta}_0 - \widehat{\boldsymbol{\beta}}^*, \boldsymbol{\xi} -\widehat{\boldsymbol{\xi}}^* ) \|^2 \ge \operatorname*{\normalfont\textrm{vec}}(\boldsymbol{\beta}_0 - \widehat{\boldsymbol{\beta}}^*, \boldsymbol{\xi} -\widehat{\boldsymbol{\xi}}^* )^\top \pmb{\Omega}_k \operatorname*{\normalfont\textrm{vec}}(\boldsymbol{\beta}_0 - \widehat{\boldsymbol{\beta}}^*, \boldsymbol{\xi} -\widehat{\boldsymbol{\xi}}^* ) \nonumber \\ &&- \frac{1}{n} \left\| [\mathbf{X}, \mathbf{V} ]^\top [\mathbf{X}, \mathbf{V} ] - \pmb{\Omega}_k\right\|_{\max} \|\operatorname*{\normalfont\textrm{vec}}(\boldsymbol{\beta}_0 - \widehat{\boldsymbol{\beta}}^*, \boldsymbol{\xi} -\widehat{\boldsymbol{\xi}}^* )\|_1^2 \nonumber\\ && \geq \operatorname*{\normalfont\textrm{vec}}(\boldsymbol{\beta}_0 - \widehat{\boldsymbol{\beta}}^*, \boldsymbol{\xi} -\widehat{\boldsymbol{\xi}}^* )^\top \pmb{\Omega}_k \operatorname*{\normalfont\textrm{vec}}(\boldsymbol{\beta}_0 - \widehat{\boldsymbol{\beta}}^*, \boldsymbol{\xi} -\widehat{\boldsymbol{\xi}}^* ) \nonumber\\ &&- 4 k_0\frac{1}{n}\left\| [\mathbf{X}, \mathbf{V}]^\top [\mathbf{X}, \mathbf{V}] - \pmb{\Omega}_k\right\|_{\max}\cdot\|\operatorname*{\normalfont\textrm{vec}}(\boldsymbol{\beta}_0 - \widehat{\boldsymbol{\beta}}^*, \boldsymbol{\xi} -\widehat{\boldsymbol{\xi}}^* )\|^2\ \nonumber \\ &&\ge \frac{1}{2}\operatorname*{\normalfont\textrm{vec}}(\boldsymbol{\beta}_0 - \widehat{\boldsymbol{\beta}}^*, \boldsymbol{\xi} -\widehat{\boldsymbol{\xi}}^* )^\top \pmb{\Omega}_k \operatorname*{\normalfont\textrm{vec}}(\boldsymbol{\beta}_0 - \widehat{\boldsymbol{\beta}}^*, \boldsymbol{\xi} -\widehat{\boldsymbol{\xi}}^* ), \end{eqnarray} where $\|\cdot \|_{\max}$ stands for the max norm for a matrix, the second inequality follows from (ref), and the third inequality follows from a development similar for (ref) and $k_0 \frac{\sqrt{\log k}}{\sqrt{n}}\to 0$. Again as explained in the beginning of this proof, we obtain \begin{eqnarray*} && \frac{1}{2}\alpha \|\operatorname*{\normalfont\textrm{vec}}(\boldsymbol{\beta}_0 - \widehat{\boldsymbol{\beta}}^*, \boldsymbol{\xi} -\widehat{\boldsymbol{\xi}}^* ) \|^2 \le \frac{3}{2}\lambda\|\operatorname*{\normalfont\textrm{vec}}(\boldsymbol{\beta}_0 - \widehat{\boldsymbol{\beta}}^*, \mathbf{S}_{\mathscr{C}} ( \boldsymbol{\xi} -\widehat{\boldsymbol{\xi}}^*) )\|_1\notag \\ &&\le \frac{3}{2}\lambda \sqrt{k_0}\|\operatorname*{\normalfont\textrm{vec}}(\boldsymbol{\beta}_0 - \widehat{\boldsymbol{\beta}}^*, \mathbf{S}_{\mathscr{C}} ( \boldsymbol{\xi} -\widehat{\boldsymbol{\xi}}^*) ) \| \le \frac{3}{2}\lambda \sqrt{k_0}\|\operatorname*{\normalfont\textrm{vec}}(\boldsymbol{\beta}_0 - \widehat{\boldsymbol{\beta}}^*, \boldsymbol{\xi} -\widehat{\boldsymbol{\xi}}^*) \| \end{eqnarray*} where the first inequality follows from (ref). We can conclude that with probability approaching 1, $\|\operatorname*{\normalfont\textrm{vec}}(\boldsymbol{\beta}_0 - \widehat{\boldsymbol{\beta}}^*, \boldsymbol{\xi} -\widehat{\boldsymbol{\xi}}^* ) \| \le \frac{3\lambda \sqrt{k_0}}{\alpha}.$ Additionally, we can obtain \begin{eqnarray*} \|\operatorname*{\normalfont\textrm{vec}}(\boldsymbol{\beta}_0 - \widehat{\boldsymbol{\beta}}^*, \boldsymbol{\xi} -\widehat{\boldsymbol{\xi}}^* )\|_1\le 4 \sqrt{k_0}\|\operatorname*{\normalfont\textrm{vec}}(\boldsymbol{\beta}_0 - \widehat{\boldsymbol{\beta}}^*, \mathbf{S}_{\mathscr{C}} (\boldsymbol{\xi} -\widehat{\boldsymbol{\xi}}^*))\|\le \frac{12\lambda k_0}{\alpha} \end{eqnarray*} with probability approaching 1, where the first inequality follows from (ref), and the second inequality follows from $\|\operatorname*{\normalfont\textrm{vec}}(\boldsymbol{\beta}_0 - \widehat{\boldsymbol{\beta}}^*, \boldsymbol{\xi}-\widehat{\boldsymbol{\xi}}^* ) \| \le \frac{3\lambda \sqrt{k_0}}{\alpha}$. The proof for (i)-(ii) is completed. (iii). To proceed, we still use the model defined in (ref), and define a new objective function here. \begin{eqnarray*} Q(\boldsymbol{\beta}, \boldsymbol{\xi}_k) =\frac{1}{n}\|\mathbf{y}- \mathbf{X}\boldsymbol{\beta} - \mathbf{V}\boldsymbol{\xi}_k\|^2 + \lambda \|\boldsymbol{\xi}_k\circ\pmb{\zeta}\|_1, \end{eqnarray*} where $\pmb{\zeta} =(\zeta_1,\ldots, \zeta_k)^\top$. Accordingly, Step 2 can be written as $(\widehat{\pmb{\beta}}^\dag, \widehat{\pmb{\xi}}_k^\dag) =\operatorname*{\arg\!\min} Q(\boldsymbol{\beta}, \boldsymbol{\xi}_k)$, where $\widehat{\pmb{\xi}}_k^\dag =(\widehat{\xi}_{k,1}^\dag,\ldots, \widehat{\xi}_{k,k}^\dag)^\top$. Due to the property of convex optimization, $(\widehat{\pmb{\beta}}^\dag, \widehat{\pmb{\xi}}_k^\dag)$ is a solution if and only if there exits a subgradient \begin{eqnarray} \mathbf{g} &\in &\{\operatorname*{\normalfont\textrm{vec}}(\mathbf{0}_{d\times 1},\mathbf{z}) \mid \mathbf{z}\in \mathbb{R}^{k},\, \notag \\ &&z_\ell = \operatorname*{\normalfont\textrm{sgn}}(\widehat{\xi}_{k,\ell}^\dag) \zeta_\ell \text{ for }\widehat{\xi}_{k,\ell}^\dag\ne 0, \text{ and } |z_\ell|\le \zeta_\ell \text{ elsewhere, for }\ell\in [k] \} \end{eqnarray} such that \begin{eqnarray} \mathbf{0} = \frac{1}{n} (\mathbf{X}, \mathbf{V})^\top (\mathbf{X}, \mathbf{V})\operatorname*{\normalfont\textrm{vec}}(\widehat{\pmb{\beta}}^\dag, \widehat{\pmb{\xi}}_k^\dag) - \frac{1}{n} (\mathbf{X}, \mathbf{V})^\top \mathbf{y}+ \frac{\lambda}{2} \mathbf{g}. \end{eqnarray} We now partition $\mathbf{V}$, $\mathbf{g}$, $\pmb{\xi}$ and $\widehat{\pmb{\xi}}_k^\dag$ based on the sets $\mathscr{C}$ and $\bar{\mathscr{C}}$. Specifically, let $\mathbf{V}_{\mathscr{C}}$, $\mathbf{g}_{\mathscr{C}}$, $\widehat{\pmb{\xi}}_{k,\mathscr{C}}^\dag$ and $\pmb{\xi}_{\mathscr{C}}$ all include the columns/elements corresponding to the set $\mathscr{C}$, and let $\mathbf{V}_{\bar{\mathscr{C}}}$ and $\mathbf{g}_{\bar{\mathscr{C}}}$ include the columns/elements corresponding to the set $\bar{\mathscr{C}}$. Here we let $\mathbf{g}_{\mathscr{C}}$ include the first $d$ 0's in $\mathbf{g}$. We now proceed. By (ref), $\operatorname*{\normalfont\textrm{sgn}}(\widehat{\pmb{\xi}}_k^\dag)=\operatorname*{\normalfont\textrm{sgn}}(\pmb{\xi})$ holds if and only if the following Karush-Kuhn-Tucker conditions hold: \begin{eqnarray} \mathbf{0}&=& \frac{1}{n}(\mathbf{X}, \mathbf{V}_{\mathscr{C}})^\top (\mathbf{X}, \mathbf{V}_{\mathscr{C}}) \operatorname*{\normalfont\textrm{vec}}(\widehat{\pmb{\beta}}^\dag -\pmb{\beta}_0, \widehat{\pmb{\xi}}_{k,\mathscr{C}}^\dag-\pmb{\xi}_{\mathscr{C}})\notag\\ &&-\frac{1}{n} (\mathbf{X}, \mathbf{V}_{\mathscr{C}})^\top (\pmb{\delta}+\mathbf{e}) +\frac{\lambda}{2}\mathbf{g}_{\mathscr{C}}, \\ \mathbf{0}&=& \frac{1}{n}\mathbf{V}_{\bar{\mathscr{C}}}^\top (\mathbf{X}, \mathbf{V}_{\mathscr{C}}) \operatorname*{\normalfont\textrm{vec}}(\widehat{\pmb{\beta}}^\dag -\pmb{\beta}_0, \widehat{\pmb{\xi}}_{k,\mathscr{C}}^\dag-\pmb{\xi}_{\mathscr{C}}) -\frac{1}{n} \mathbf{V}_{\bar{\mathscr{C}}}^\top (\pmb{\delta}+\mathbf{e}) +\frac{\lambda}{2}\mathbf{g}_{\bar{\mathscr{C}}}, \end{eqnarray} and \begin{eqnarray} \operatorname*{\normalfont\textrm{sgn}}(\xi_{\ell})(\widehat{\xi}_{k, \ell}^\dag -\xi_{\ell})>-|\xi_{\ell}|\quad\text{for}\quad\ell \in \mathscr{C}. \end{eqnarray} We justify (ref) and (ref) first. Define $\pmb{\Omega}_{\mathscr{C}} =\begin{pmatrix} \mathbf{A}_{11} & \mathbf{A}_{12} \\ \mathbf{A}_{12}^\top & \mathbf{I}_k \\ \end{pmatrix}$ and $\mathbf{A}_{12}\coloneqq (\pmb{\eta}_1,\cdots , \pmb{\eta}_{k_0})$, where $\mathbf{A}_{11}=E[\mathbf{x}\mathbf{x}^\top]$ and $\pmb{\eta}_\ell=E[\mathbf{x}\psi_\ell(\mathbf{x})]$. By the development of $\pmb{\Omega}_k$, it is easy to know that $\pmb{\Omega}_{\mathscr{C}}$ is invertible and $\rho_{\min}(\pmb{\Omega}_{\mathscr{C}})>0$. Note that we have shown that \begin{eqnarray*} &&\left|\rho_{\min}(\frac{1}{n}(\mathbf{X}, \mathbf{V}_{\mathscr{C}})^\top (\mathbf{X}, \mathbf{V}_{\mathscr{C}})) -\rho_{\min}(\pmb{\Omega}_{\mathscr{C}} )\right| \notag \\ &\le & (d+k_0) \left\| \frac{1}{n}(\mathbf{X}, \mathbf{V}_{\mathscr{C}})^\top (\mathbf{X}, \mathbf{V}_{\mathscr{C}})- \pmb{\Omega}_{\mathscr{C}} \right\|_{\max}\lesssim k_0\frac{\sqrt{\log k}}{\sqrt{n}} \asymp k_0\lambda, \end{eqnarray*} where $\rho_{\min}(\cdot)$ stands for the minimum eigenvalue, the second step follow from the development of (ref), and the last step follows from $\lambda\asymp \frac{\sqrt{\log k}}{\sqrt{n}} $. Thus, we can write \begin{eqnarray} &&\left\| \left(\frac{1}{n}(\mathbf{X}, \mathbf{V}_{\mathscr{C}})^\top (\mathbf{X}, \mathbf{V}_{\mathscr{C}})\right)^{-1} \right\|_\infty\le \sqrt{k_0+d}\left\| \left(\frac{1}{n}(\mathbf{X}, \mathbf{V}_{\mathscr{C}})^\top (\mathbf{X}, \mathbf{V}_{\mathscr{C}})\right)^{-1}\right\|_2\notag \\ &\le & \sqrt{k_0+d}/\rho_{\min}(\frac{1}{n}(\mathbf{X}, \mathbf{V}_{\mathscr{C}})^\top (\mathbf{X}, \mathbf{V}_{\mathscr{C}})) \lesssim \sqrt{k_0}. \end{eqnarray} From (ref), we obtain the following expression \begin{eqnarray} &&\operatorname*{\normalfont\textrm{vec}}(\widehat{\pmb{\beta}}^\dag -\pmb{\beta}_0, \widehat{\pmb{\xi}}_{k,\mathscr{C}}^\dag-\pmb{\xi}_{\mathscr{C}})\notag \\ &=&\left(\frac{1}{n}(\mathbf{X}, \mathbf{V}_{\mathscr{C}})^\top (\mathbf{X}, \mathbf{V}_{\mathscr{C}})\right)^{-1}\left(\frac{1}{n} (\mathbf{X}, \mathbf{V}_{\mathscr{C}})^\top (\pmb{\delta}+\mathbf{e}) -\frac{\lambda}{2}\mathbf{g}_{\mathscr{C}}\right). \end{eqnarray} The right hand side of (ref) can be further written as follows: \begin{eqnarray} &&\left\|\left(\frac{1}{n}(\mathbf{X}, \mathbf{V}_{\mathscr{C}})^\top (\mathbf{X}, \mathbf{V}_{\mathscr{C}})\right)^{-1}\left(\frac{1}{n} (\mathbf{X}, \mathbf{V}_{\mathscr{C}})^\top (\pmb{\delta}+\mathbf{e}) -\frac{\lambda}{2}\mathbf{g}_{\mathscr{C}}\right)\right\|_\infty\notag \\ &\lesssim &\left\| \left(\frac{1}{n}(\mathbf{X}, \mathbf{V}_{\mathscr{C}})^\top (\mathbf{X}, \mathbf{V}_{\mathscr{C}})\right)^{-1} \right\|_\infty\left(\left\|\frac{1}{n} (\mathbf{X}, \mathbf{V}_{\mathscr{C}})^\top (\pmb{\delta}+\mathbf{e}) \right\|_\infty + \frac{\lambda}{2}\max_{\ell\in \mathscr{C}} \zeta_{\ell}\right)\notag \\ &\lesssim & \sqrt{k_0}\left(\frac{\sqrt{\log k}}{\sqrt{n}} + \frac{\lambda}{2}\max_{\ell\in \mathscr{C}} \zeta_{\ell}\right) \lesssim \lambda\sqrt{k_0} (1+ \max_{\ell\in \mathscr{C}} \zeta_{\ell} ), \end{eqnarray} where the second step follows from (ref) and (ref). Thus, (ref) is satisfied provided $\min_{\ell \in \mathscr{C}} |\xi_{\ell}| \gg \lambda\sqrt{k_0} (1+ \max_{\ell\in \mathscr{C}} \zeta_{\ell} ).$ We now turn to (ref). Combining (ref) and (ref), we shall show that for $\forall\ell \in \bar{\mathscr{C}}$ \begin{eqnarray} \frac{\lambda}{2}\zeta_{\ell} &\ge &\left|\frac{1}{n}\mathbf{V}_{\ell}^\top (\mathbf{X}, \mathbf{V}_{\mathscr{C}}) \left(\frac{1}{n}(\mathbf{X}, \mathbf{V}_{\mathscr{C}})^\top (\mathbf{X}, \mathbf{V}_{\mathscr{C}})\right)^{-1}\left(\frac{1}{n} (\mathbf{X}, \mathbf{V}_{\mathscr{C}})^\top (\pmb{\delta}+\mathbf{e}) -\frac{\lambda}{2}\mathbf{g}_{\mathscr{C}}\right) \right| \notag \\ &&+\left| \frac{1}{n} \mathbf{V}_{\ell}^\top (\pmb{\delta}+\mathbf{e})\right|, \end{eqnarray} where $\mathbf{V}_{\ell}$ stands for the $\ell^{th}$ column of $\mathbf{V}$. Note that for $\forall\ell \in \bar{\mathscr{C}}$ \begin{eqnarray*} &&\left|\frac{1}{n}\mathbf{V}_{ \ell}^\top (\mathbf{X}, \mathbf{V}_{\mathscr{C}}) \left(\frac{1}{n}(\mathbf{X}, \mathbf{V}_{\mathscr{C}})^\top (\mathbf{X}, \mathbf{V}_{\mathscr{C}})\right)^{-1}\left(\frac{1}{n} (\mathbf{X}, \mathbf{V}_{\mathscr{C}})^\top (\pmb{\delta}+\mathbf{e}) -\frac{\lambda}{2}\mathbf{g}_{\mathscr{C}}\right) \right| \notag \\ &\le &\sqrt{k_0}\left\| \frac{1}{n}(\mathbf{X}, \mathbf{V}_{\mathscr{C}})^\top (\mathbf{X}, \mathbf{V}_{\mathscr{C}}) \right\|_{\max} \notag \\ &&\cdot \left\| \left(\frac{1}{n}(\mathbf{X}, \mathbf{V}_{\mathscr{C}})^\top (\mathbf{X}, \mathbf{V}_{\mathscr{C}})\right)^{-1}\left(\frac{1}{n} (\mathbf{X}, \mathbf{V}_{\mathscr{C}})^\top (\pmb{\delta}+\mathbf{e}) -\frac{\lambda}{2}\mathbf{g}_{\mathscr{C}}\right) \right\|_{\infty}\notag \\ &\lesssim & \sqrt{k_0}\|\operatorname*{\normalfont\textrm{diag}}\{E[\mathbf{x}_i\mathbf{x}_i^\top], \mathbf{I}_{k_0} \}\|_{\max}\lambda\sqrt{k_0} (1+ \max_{\ell\in \mathscr{C}} \zeta_{\ell} )\lesssim \lambda k_0 (1+ \max_{\ell\in \mathscr{C}} \zeta_{\ell} ) \end{eqnarray*} and $\left| \frac{1}{n} \mathbf{V}_{\ell}^\top (\pmb{\delta}+\mathbf{e})\right|\lesssim \frac{\sqrt{\log(k)}}{\sqrt{n}}\asymp \lambda$. Thus, to show (ref), we require $\min_{\ell \in \bar{\mathscr{C}}} \zeta_{\ell}\gg k_0 \left(1+ \max_{\ell\in \mathscr{C}} \zeta_{\ell}\right)$. Putting everything together, the proof of (iii) is now completed. (iv). Note that (ref) requires \begin{eqnarray*} \mathbf{0}&=& \frac{1}{n}(\mathbf{X}, \mathbf{V}_{\mathscr{C}})^\top (\mathbf{X}, \mathbf{V}_{\mathscr{C}}) \operatorname*{\normalfont\textrm{vec}}(\widehat{\pmb{\beta}}^\dag -\pmb{\beta}_0, \widehat{\pmb{\xi}}_{k,\mathscr{C}}^\dag-\pmb{\xi}_{\mathscr{C}})-\frac{1}{n} (\mathbf{X}, \mathbf{V}_{\mathscr{C}})^\top (\pmb{\delta}+\mathbf{e}) +\frac{\lambda}{2}\mathbf{g}_{\mathscr{C}}, \end{eqnarray*} wherein the first $d$ rows are \begin{eqnarray*} \mathbf{0}&=& \frac{1}{n}\mathbf{X}^\top (\mathbf{X}, \mathbf{V}_{\mathscr{C}}) \operatorname*{\normalfont\textrm{vec}}(\widehat{\pmb{\beta}}^\dag -\pmb{\beta}_0, \widehat{\pmb{\xi}}_{k,\mathscr{C}}^\dag-\pmb{\xi}_{\mathscr{C}}) -\frac{1}{n} \mathbf{X}^\top (\pmb{\delta}+\mathbf{e}) \end{eqnarray*} due to the fact that the first $d$ elements of $\mathbf{g}_{\mathscr{C}}$ are 0 by the design. Here $\widehat{\pmb{\beta}}^\dag$ is from the first order condition when minimizing $Q(\pmb{\beta},\pmb{\xi}_{k,\mathscr{C}}^\dag) =\| \mathbf{y}-\mathbf{X}\boldsymbol{\beta} -\mathbf{V}_{\mathscr{C}}\pmb{\xi}_{k,\mathscr{C}}^\dag\|^2 $. Note also that \begin{eqnarray*} &&Q(\pmb{\beta},\pmb{\xi}_{k,\mathscr{C}}^\dag) =\| \mathbf{y}-\mathbf{X}\boldsymbol{\beta} -\mathbf{V}_{\mathscr{C}}\pmb{\xi}_{k,\mathscr{C}}^\dag\|^2 \notag \\ &=&\|\mathbf{M}_{\mathbf{V}_{\mathscr{C}}}(\mathbf{y}- \mathbf{X}\boldsymbol{\beta})\|^2+ \|\mathbf{P}_{\mathbf{V}_{\mathscr{C}}}(\mathbf{y}- \mathbf{X}\boldsymbol{\beta}) -\mathbf{V}_{\mathscr{C}}\pmb{\xi}_{k,\mathscr{C}}^\dag\|^2\notag \\ &=&\|\mathbf{M}_{\mathbf{V}_{\mathscr{C}}}(\mathbf{y}- \mathbf{X}\boldsymbol{\beta})\|^2+ \|\mathbf{V}_{\mathscr{C}}(\frac{1}{n}\mathbf{V}_{\mathscr{C}}^\top\mathbf{y}-\frac{1}{n} \mathbf{V}_{\mathscr{C}}^\top\mathbf{X}(\boldsymbol{\beta}+\Delta\boldsymbol{\beta}) - \pmb{\xi}_{k,\mathscr{C}}^\dag +\frac{1}{n} \mathbf{V}_{\mathscr{C}}^\top\mathbf{X} \Delta\boldsymbol{\beta})\|^2\notag \\ &\eqqcolon&\|\mathbf{M}_{\mathbf{V}_{\mathscr{C}}}(\mathbf{y}- \mathbf{X}\boldsymbol{\beta})\|^2+ \|\mathbf{V}_{\mathscr{C}}(\frac{1}{n}\mathbf{V}_{\mathscr{C}}^\top\mathbf{y}-\frac{1}{n} \mathbf{V}_{\mathscr{C}}^\top\mathbf{X}\widetilde{\boldsymbol{\beta}} - \widetilde{\pmb{\xi}}_{k,\mathscr{C}}^\dag )\|^2 \end{eqnarray*} wherein $\widetilde{\boldsymbol{\beta}} = \boldsymbol{\beta}+\Delta\boldsymbol{\beta}$, $\widetilde{\pmb{\xi}}_{k,\mathscr{C}}^\dag = \pmb{\xi}_{k,\mathscr{C}}^\dag -\frac{1}{n} \mathbf{V}_{\mathscr{C}}^\top\mathbf{X} \Delta\boldsymbol{\beta}$, and $\Delta\boldsymbol{\beta}$ stands for a generic vector. Note that we partition $Q(\pmb{\beta},\pmb{\xi}_{k,\mathscr{C}}^\dag)$ into two parts formed by two orthogonal spaces, of which both can reach its own minimum. The first term on the right hand side is uniquely minimized at $\widehat{\pmb{\beta}}^\dag =(\mathbf{X}^\top\mathbf{M}_{\mathbf{V}_{\mathscr{C}}}\mathbf{X})^{-1} \mathbf{X}^\top\mathbf{M}_{\mathbf{V}_{\mathscr{C}}} \mathbf{y}$, and the second term on the right hand side researches its minimum but the value of $\widetilde{\boldsymbol{\beta}}$ and $\widetilde{\pmb{\xi}}_{k,\mathscr{C}}^\dag$ may move jointly. Thus, in order to minimize $Q(\pmb{\beta},\pmb{\xi}_{k,\mathscr{C}}^\dag)$, one should get $\widehat{\pmb{\xi}}_{k,\mathscr{C}}^\dag$ first, and then decide $\widehat{\pmb{\beta}}^\dag$ accordingly. From here, the CLT follows from the proof of Theorem 3.1 immediately.

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